--- license: openrail++ library_name: diffusers tags: - text-to-image - text-to-image - diffusers-training - diffusers - lora - template:sd-lora - stable-diffusion-xl - stable-diffusion-xl-diffusers base_model: SDXL_model instance_prompt: a photo of pll cat widget: - text: a photo of pll cat by the sea output: url: image_0.png - text: a photo of pll cat by the sea output: url: image_1.png - text: a photo of pll cat by the sea output: url: image_2.png - text: a photo of pll cat by the sea output: url: image_3.png --- # SDXL LoRA DreamBooth - YongjieNiu/prior_lora-pll-cat-100-1-500 ## Model description These are YongjieNiu/prior_lora-pll-cat-100-1-500 LoRA adaption weights for SDXL_model. The weights were trained using [DreamBooth](https://dreambooth.github.io/). LoRA for the text encoder was enabled: False. Special VAE used for training: VAE. ## Trigger words You should use a photo of pll cat to trigger the image generation. ## Download model Weights for this model are available in Safetensors format. [Download](YongjieNiu/prior_lora-pll-cat-100-1-500/tree/main) them in the Files & versions tab. ## Intended uses & limitations #### How to use ```python # TODO: add an example code snippet for running this diffusion pipeline ``` #### Limitations and bias [TODO: provide examples of latent issues and potential remediations] ## Training details [TODO: describe the data used to train the model]