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The Gall–Peters projection is a rectangular map projection that maps all areas such that they have the correct sizes relative to each other. Like any equal-area projection, it achieves this goal by distorting most shapes. The projection is a particular example of the cylindrical equal-area projection with latitudes 45° north and south as the regions on the map that have no distortion.
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The projection is named after James Gall and Arno Peters. Gall is credited with describing the projection in 1855 at a science convention. He published a paper on it in 1885.[1] Peters brought the projection to a wider audience beginning in the early 1970s by means of the "Peters World Map". The name "Gall–Peters projection" seems to have been used first by Arthur H. Robinson in a pamphlet put out by the American Cartographic Association in 1986.[2]
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Maps based on the projection are promoted by UNESCO, and they are also widely used by British schools.[3] The U.S. state of Massachusetts and Boston Public Schools began phasing in these maps in March 2017, becoming the first public school district and state in the United States to adopt Gall–Peters maps as their standard.[4]
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The Gall–Peters projection achieved notoriety in the late 20th century as the centerpiece of a controversy about the political implications of map design.[5]
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The projection is conventionally defined as:
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where λ is the longitude from the central meridian in degrees, φ is the latitude, and R is the radius of the globe used as the model of the earth for projection. For longitude given in radians, remove the π/180° factors.
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Stripping out unit conversion and uniform scaling, the formulae may be written:
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where λ is the longitude from the central meridian (in radians), φ is the latitude, and R is the radius of the globe used as the model of the earth for projection. Hence the sphere is mapped onto the vertical cylinder, and the cylinder is stretched to double its length. The stretch factor, 2 in this case, is what distinguishes the variations of cylindric equal-area projection.
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The various specializations of the cylindric equal-area projection differ only in the ratio of the vertical to horizontal axis. This ratio determines the standard parallel of the projection, which is the parallel at which there is no distortion and along which distances match the stated scale. There are always two standard parallels on the cylindric equal-area projection, each at the same distance north and south of the equator. The standard parallels of the Gall–Peters are 45° N and 45° S. Several other specializations of the equal-area cylindric have been described, promoted, or otherwise named.[6][7][8]
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The Gall–Peters projection was first described in 1855 by clergyman James Gall, who presented it along with two other projections at the Glasgow meeting of the British Association for the Advancement of Science (the BA). He gave it the name "orthographic" and formally published his work in 1885 in the Scottish Geographical Magazine.[1] The projection is suggestive of the Orthographic projection in that distances between parallels of the Gall–Peters are a constant multiple of the distances between the parallels of the orthographic. That constant is √2.
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The name "Gall–Peters projection" seems to have been used first by Arthur H. Robinson in a pamphlet put out by the American Cartographic Association in 1986.[2] Before 1973 it had been known, when referred to at all, as the "Gall orthographic" or "Gall's orthographic." Most Peters supporters refer to it only as the "Peters projection." During the years of controversy the cartographic literature tended to mention both attributions, settling on one or the other for the purposes of the article. In recent years "Gall–Peters" seems to dominate.
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In 1967, Arno Peters, a German filmmaker, devised a map projection identical to Gall's orthographic projection and presented it in 1973 as a "new invention". He promoted it as a superior alternative to the Mercator projection, which was suited to navigation but also used commonly in world maps. The Mercator projection increasingly inflates the sizes of regions according to their distance from the equator. This inflation results, for example, in a representation of Greenland that is larger than Africa, which has a geographic area 14 times greater than Greenland's. Since much of the technologically underdeveloped world lies near the equator, these countries appear smaller on a Mercator and therefore, according to Peters, seem less significant.[7]:155 On Peters's projection, by contrast, areas of equal size on the globe are also equally sized on the map. By using his "new" projection, Peters argued that poorer, less powerful nations could be restored to their rightful proportions. This reasoning has been picked up by many educational and religious bodies, leading to adoption of the Gall–Peters projection among some socially concerned groups, including Oxfam,[9] National Council of Churches,[10] New Internationalist magazine,[11] and the Mennonite Central Committee.[12] However, Peters's choice of 45° N/S for the standard parallels means that the regions displayed with highest accuracy include Europe and the US, and not the tropics.
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Peters's original description of the projection for his map contained a geometric error that, taken literally, implies standard parallels of 46°02′ N/S. However the text accompanying the description made it clear that he had intended the standard parallels to be 45° N/S, making his projection identical to Gall's orthographic.[13] In any case, the difference is negligible in a world map.
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At first, the cartographic community largely ignored Peters's foray into cartography. The preceding century had already witnessed many campaigns for new projections with little visible result. Just twenty years earlier, for example, Trystan Edwards described and promoted his own eponymous projection, disparaging the Mercator, and recommending his projection as the solution.[14] Peters's projection differed from Edwards's only in height-to-width ratio. More problematic, Peters's projection was identical to one that was already over a century old, though he probably did not realize it.[7] That projection—Gall's orthographic—passed unnoticed when it was announced in 1855.
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Beyond the lack of novelty in the projection itself, the claims Peters made about the projection were also familiar to cartographers. Just as in the case of Peters, earlier projections generally were promoted as alternatives to the Mercator. Inappropriate use of the Mercator projection in world maps and the size disparities figuring prominently in Peters's arguments against the Mercator projection had been remarked upon for centuries and quite commonly in the 20th century.[14][15][16][17][18][19][20] As early as 1943, Stewart notes this phenomenon and compares the quest for the perfect projection to "squaring the circle or making pi come out even"[21] because the mathematics that governs map projections just does not permit development of a map projection that is objectively significantly better than the hundreds already devised. Even Peters's politicized interpretation of the common use of Mercator was nothing new, with Kelloway's 1946 text mentioning a similar controversy.[17]
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Cartographers had long despaired over publishers' inapt use of the Mercator.[22][23][24][25] A 1943 New York Times editorial stated that "The time has come to discard [the Mercator] for something that represents the continents and directions less deceptively ... Although its usage ... has diminished ... it is still highly popular as a wall map apparently in part because, as a rectangular map, it fills a rectangular wall space with more map, and clearly because its familiarity breeds more popularity."[26] Because of the lack of novelty both in the projection Peters devised and in the rhetoric surrounding its promotion, the cartographic community had no reason to think Peters would succeed any more than Edwards or his predecessors had.[26]:165
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Peters, however, launched his campaign in a different world from that of Edwards. He announced his map at a time when themes of social justice resonated strongly in academia and politics. Suggesting "cartographic imperialism", Peters found ready audiences. The campaign was bolstered by the claim that the Peters projection was the only "area-correct" map. Other claims included "absolute angle conformality", "no extreme distortions of form", and "totally distance-factual".[27]
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All of those claims were erroneous.[28] Some of the oldest projections are equal-area (the sinusoidal projection is also known as the "Mercator equal-area projection"), and hundreds have been described, refuting any implication that Peters's map is special in that regard. In any case, Mercator was not the pervasive projection Peters made it out to be: a wide variety of projections has always been used in world maps.[29] Peters's chosen projection suffers extreme distortion in the polar regions, as any cylindrical projection must, and its distortion along the equator is considerable. Several scholars have remarked on the irony of the projection's undistorted presentation of the mid latitudes, including Peters's native Germany, at the expense of the low latitudes, which host more of the technologically underdeveloped nations.[30][31] The claim of distance fidelity is particularly problematic: Peters's map lacks distance fidelity everywhere except along the 45th parallels north and south, and then only in the direction of those parallels. No world projection is good at preserving distances everywhere; Peters's and all other cylindric projections are especially bad in that regard because east-west distances inevitably balloon toward the poles.[28][32]
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The cartographic community met Peters's 1973 press conference with amusement and mild exasperation, but little activity beyond a few articles commenting on the technical aspects of Peters's claims. In the ensuing years, however, it became clear that Peters and his map were no flash in the pan. By 1980 many cartographers had turned overtly hostile to his claims. In particular, Peters writes in The New Cartography,
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Philosophers, astronomers, historians, popes and mathematicians have all drawn global maps long before cartographers as such existed. Cartographers appeared in the "Age of Discovery", which developed into the Age of European Conquest and Exploitation and took over the task of making maps.
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By the authority of their profession they have hindered its development. Since Mercator produced his global map over four hundred years ago for the age of Europeans world domination, cartographers have clung to it despite its having been long outdated by events. They have sought to render it topical by cosmetic corrections.
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... The European world concept, as the last expression of a subjective global view of primitive peoples, must give way to an objective global concept.
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The cartographic profession is, by its retention of old precepts based on the Eurocentric global concept, incapable of developing this egalitarian world map which alone can demonstrate the parity of all peoples of the earth.[33]
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This attack galled the cartographic community. Their most emphatic refutation of Peters's assertions was the long list of cartographers who, over the preceding century, had formally expressed frustration at publishers' overuse of the Mercator, as noted above.[26] Many of those cartographers had already developed projections they explicitly promoted as alternatives to the Mercator, including the most influential American cartographers of the twentieth century: John Paul Goode (Goode homolosine projection), Erwin Raisz (Armadillo projection), and Arthur H. Robinson (Robinson projection). Hence the cartographic community viewed Peters's narrative as ahistorical and mean-spirited.
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The two camps never made any real attempts toward reconciliation. The Peters camp largely ignored the protests of the cartographers. Peters maintained there should be "one map for one world"[34]—his—and did not acknowledge the prior art of Gall[28] until the controversy had largely run its course, late in his life. While Peters likely reinvented the projection independently, his unscholarly conduct and refusal to engage the cartographic community undoubtedly contributed to the polarization and impasse.[5]
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Frustrated by some very visible successes and mounting publicity stirred up by the industry that had sprung up around the Peters map, the cartographic community began to plan more coordinated efforts to restore balance, as they saw it. The 1980s saw a flurry of literature directed against the Peters phenomenon. Though Peters's map was not singled out, the controversy motivated the American Cartographic Association (now Cartography and Geographic Information Society) to produce a series of booklets (including Which Map Is Best[2]) designed to educate the public about map projections and distortion in maps. In 1989 and 1990, after some internal debate, seven North American geographic organizations adopted the following resolution,[35][36] which rejected all rectangular world maps, a category that includes both the Mercator and the Gall–Peters projections:
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WHEREAS, the earth is round with a coordinate system composed entirely of circles, and
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WHEREAS, flat world maps are more useful than globe maps, but flattening the globe surface necessarily greatly changes the appearance of Earth's features and coordinate systems, and
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WHEREAS, world maps have a powerful and lasting effect on people's impressions of the shapes and sizes of lands and seas, their arrangement, and the nature of the coordinate system, and
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WHEREAS, frequently seeing a greatly distorted map tends to make it "look right",
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THEREFORE, we strongly urge book and map publishers, the media and government agencies to cease using rectangular world maps for general purposes or artistic displays. Such maps promote serious, erroneous conceptions by severely distorting large sections of the world, by showing the round Earth as having straight edges and sharp corners, by representing most distances and direct routes incorrectly, and by portraying the circular coordinate system as a squared grid. The most widely displayed rectangular world map is the Mercator (in fact a navigational diagram devised for nautical charts), but other rectangular world maps proposed as replacements for the Mercator also display a greatly distorted image of the spherical Earth.
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One map society, the North American Cartographic Information Society (NACIS), declined to endorse the 1989 resolution, although no reasons were given.[citation needed]
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The geographic and cartographic communities did not unanimously disparage the Peters World Map. Some cartographers, including J. Brian Harley, have credited the Peters phenomenon with demonstrating the social implications of map projections, at the very least.[37] Crampton sees the condemnation from the cartographic community as reactionary and perhaps demonstrative of immaturity in the profession, given that all maps are political.[5] Denis Wood sees the map as one of many useful tools.[34] Lastly, Terry Hardaker of Oxford Cartographers Limited, sympathetic to Peters's mission, became the map's official cartographer when Peters, overwhelmed by the technical aspects of cartography, sought to pass on those responsibilities.[34]
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Notes
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Further reading
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Reportedly haunted locations:
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In religion, a prophet is an individual who is regarded as being in contact with a divine being and is said to speak on that entity's behalf, serving as an intermediary with humanity by delivering messages or teachings from the supernatural source to other people.[1][2] The message that the prophet conveys is called a prophecy.
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Claims of prophethood have existed in many cultures throughout history, including Judaism, Christianity, Islam, in ancient Greek religion, Zoroastrianism, Manichaeism, and many others.
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The English word prophet is a compound Greek word, from pro (in advance) and the verb phesein (to tell); thus, a προφήτης (prophétés) is someone who foretells future events, and also conveys messages from the divine to humans; in a different interpretation, it means advocate or speaker.
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In Hebrew, the word נָבִיא (nāvî), "spokesperson", traditionally translates as "prophet".[3] The second subdivision of the Tanakh, (Nevi'im), is devoted to the Hebrew prophets. The meaning of navi is perhaps described in Deuteronomy 18:18,[4] where God said, "...and I will put My words in his mouth, and he shall speak unto them all that I shall command him." Thus, the navi was thought to be the "mouth" of God. The root nun-bet-alef ("navi") is based on the two-letter root nun-bet which denotes hollowness or openness; to receive transcendental wisdom, one must make oneself "open".[5]
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In addition to writing and speaking messages from God, Israelite or Judean nevi'im ("spokespersons", "prophets") often acted out prophetic parables in their life.[7] For example, in order to contrast the people's disobedience with the obedience of the Rechabites, God has Jeremiah invite the Rechabites to drink wine, in disobedience to their ancestor's command. The Rechabites refuse, for which God commends them.[8][9] Other prophetic parables acted out by Jeremiah include burying a linen belt so that it gets ruined to illustrate how God intends to ruin Judah's pride.[10][11][11][12] Likewise, Jeremiah buys a clay jar and smashes it in the Valley of Ben Hinnom in front of elders and priests to illustrate that God will smash the nation of Judah and the city of Judah beyond repair.[13] God instructs Jeremiah to make a yoke from wood and leather straps and to put it on his own neck to demonstrate how God will put the nation under the yoke of Nebuchadnezzar, king of Babylon.[14] In a similar way, the prophet Isaiah had to walk stripped and barefoot for three years to illustrate the coming captivity,[15] and the prophet Ezekiel had to lie on his side for 390 days and eat measured food to illustrate the coming siege.[16]
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The prophetic assignment is not always portrayed as positive in the Hebrew Bible,[17][18][19] and prophets were often the target of persecution and opposition.[20] God's personal prediction for Jeremiah, "Attack you they will, overcome you they can't,"[21] was performed many times in the biblical narrative as Jeremiah warned of destruction of those who continued to refuse repentance and accept more moderate consequences.[20][22] In return for his adherence to God's discipline and speaking God's words, Jeremiah was attacked by his own brothers,[23] beaten and put into the stocks by a priest and false prophet,[24][25] imprisoned by the king,[26] threatened with death,[27] thrown into a cistern by Judah's officials,[28] and opposed by a false prophet.[29] Likewise, Isaiah was told by his hearers who rejected his message, "Leave the way! Get off the path! Let us hear no more about the Holy One of Israel!"[18][30] The life of Moses being threatened by Pharaoh is another example.[31]
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According to I Samuel 9:9,[32] the old name for navi is ro'eh, רֹאֶה, which literally means "Seer". That could document an ancient shift, from viewing prophets as seers for hire to viewing them as moral teachers. L.C. Allen (1971) comments that in the First Temple Era, there were essentially seer-priests, who formed a guild, divined, performed rituals and sacrifices, and were scribes, and then there were canonical prophets, who did none of these (and were against divination) and had instead a message to deliver.[33] The seer-priests were usually attached to a local shrine or temple, such as Shiloh, and initiated others as priests in that priesthood: it was a mystical craft-guild with apprentices and recruitment. Canonical prophets were not organised this way.
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Some examples of prophets in the Tanakh include Abraham, Moses, Miriam, Isaiah, Samuel, Ezekiel, Malachi, and Job. In Jewish tradition Daniel is not counted in the list of prophets.
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A Jewish tradition suggests that there were twice as many prophets as the number which left Egypt, which would make 1,200,000 prophets.[34] The Talmud recognizes the existence of 48 male prophets who bequeathed permanent messages to mankind.[34] According to the Talmud there were also seven women who are counted as prophetesses whose message bears relevance for all generations: Sarah, Miriam, Devorah, Hannah (mother of the prophet Samuel), Abigail (a wife of King David), Huldah (from the time of Jeremiah), and Esther.[34] The Talmudic and Biblical commentator Rashi points out that Rebecca, Rachel, and Leah were also prophets.[35]
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Isaiah 8:3-4[36] refers he married "the prophetess", which conceived and gave to him a son, named by God Mahèr-salàl-cash-baz. Her name isn't elsewhere specified.
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Prophets in Tanakh are not always Jews.[34] The story of Balaam in Numbers 22 describes a non-Jewish prophet.[37] According to the Talmud, Obadiah is said to have been a convert to Judaism.
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The last nevi'im ("spokespersons", "prophets") mentioned in the Jewish Bible are Haggai, Zechariah, and Malachi, all of whom lived at the end of the 70-year Babylonian exile. The Talmud (Sanhedrin 11a) states that Haggai, Zachariah, and Malachi were the last prophets, and nowadays only the "Bath Kol" (בת קול, lit. daughter of a voice, "voice of God") exists.
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In Christianity, a prophet (or seer) is one inspired by God through the Holy Spirit to deliver a message. Some Christian denominations limit a prophet's message to words intended only for the entire church congregation, excluding personal messages not intended for the body of believers; but in the Bible on a number of occasions prophets were called to deliver personal messages.[38] The reception of a message is termed revelation and the delivery of the message is termed prophecy.
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The term "prophet" applies to those who receive public or private revelation. Public revelation, in Catholicism, is part of the Deposit of faith, the revelation of which was completed by Jesus; whereas private revelation does not add to the Deposit. The term "deposit of faith" refers to the entirety of Jesus Christ's revelation, and is passed to successive generations in two different forms, sacred scripture (the Bible) and sacred tradition.
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The Bible terms anyone who claims to speak God's words or to teach in his name without being a prophet a false prophet.[citation needed] One Old Testament text in Deuteronomy[39] contains a warning against those who prophesy events which do not come to pass and says they should be put to death. Elsewhere a false prophet may be someone who is purposely trying to deceive, is delusional, under the influence of Satan or is speaking from his own spirit.[40]
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Some Christians believe that the Holy Spirit gives spiritual gifts to Christians. These may include prophecy, tongues, miraculous healing ability, and discernment (Matthew 12:32 KJV "Whosoever speaketh a word against the Son of Man, it shall be forgiven him: but whosoever speaketh against the Holy Ghost, it shall not be forgiven him, neither in this world, neither in the world to come."). Cessationists believe that these gifts were given only in New Testament times and that they ceased after the last apostle died.
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New Testament passages that explicitly discuss prophets existing after the death and resurrection of Christ include Revelation 11:10,[41] Matthew 10:40–41 and 23:34,[42] John 13:20 and 15:20[43] and Acts 11:25–30, 13:1 and 15:32.[44]
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The Didache gives extensive instruction in how to distinguish between true and false prophets, as well as commands regarding tithes to prophets in the church.[45] Irenaeus, wrote of 2nd-century believers with the gift of prophecy,[46] while Justin Martyr argued in his Dialogue with Trypho that prophets were not found among the Jews in his time, but that the church had prophets.[47] The Shepherd of Hermas describes revelation in a vision regarding the proper operation of prophecy in the church.[48] Eusebius mentions that Quadratus and Ammia of Philadelphia were both prominent prophets following the age of the Twelve Apostles.[49][50] Tertullian, writing of the church meetings of the Montanists (to whom he belonged), described in detail the practice of prophecy in the 2nd-century church.[51]
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A number of later Christian saints were claimed to have powers of prophecy, such as Columba of Iona (521-597), Saint Malachy (1094-1148) or Padre Pio (1887-1968).[52] Marian apparitions like those at Fatima in 1917 or at Kibeho in Rwanda in the 1980s often included prophetic predictions regarding the future of the world as well as of the local areas they occurred in.[53]
|
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|
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Prophetic movements in particular can be traced throughout the Christian Church's history, expressing themselves in (for example) Montanism, Novatianism, Donatism, Franciscanism, Anabaptism, Camisard enthusiasm, Puritanism, Quakerism, Quietism, Lutheranism[54] and Pietism. Modern Pentecostals and Charismatics, members of movements which together comprised approximately 584 million people as of 2011[update],[55] believe in the contemporary function of the gift of prophecy, and some in these movements allow for idea that God may continue to gift the church with some individuals who are prophets.
|
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Some Christian sects recognize the existence of a "modern-day" prophets. One such denomination is The Church of Jesus Christ of Latter-day Saints, which teaches that God still communicates with mankind through prophecy.
|
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The Quran identifies a number of men as "Prophets of Islam" (Arabic: نبي nabī; pl. أنبياء anbiyāʾ). Muslims believe such individuals were assigned a special mission by God to guide humanity. Besides Muhammad, this includes prophets such as Abraham (Ibrāhīm), Moses (Mūsā) and Jesus (ʿĪsā).
|
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Although only twenty-five prophets[56] are mentioned by name in the Quran, a hadith (no. 21257 in Musnad Ahmad ibn Hanbal)[57] mentions that there were (more or less) 124,000 prophets in total throughout history. Other traditions place the number of prophets at 224,000. Some scholars hold that there are an even greater number in the history of mankind, and only God knows. The Quran says that God has sent a prophet to every group of people throughout time, and that Muhammad is the last of the prophets, sent for the whole of humankind.[58] The message of all the prophets is believed to be the same. In Islam, all prophetic messengers are prophets (such as Adam, Noah, Abraham, Moses, Jesus, and Muhammad) though not all prophets are prophetic messengers. The primary distinction is that a prophet is required to demonstrate God's law through his actions, character, and behavior without necessarily calling people to follow him, while a prophetic messenger is required to pronounce God's law (i.e. revelation) and call his people to submit and follow him. Muhammad is distinguished from the rest of the prophetic messengers and prophets in that he was commissioned by God to be the prophetic messenger to all of mankind. Many of these prophets are also found in the texts of Judaism (The Torah, the Prophets, and the Writings) and Christianity.[59]
|
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|
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Muslims often refer to Muhammad as "the Prophet", in the form of a noun.[60][61][62][63] Jesus is the result of a virgin birth in Islam as in Christianity, and is regarded as a prophet.[64]
|
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Traditionally, four prophets are believed to have been sent holy books: the Torah (Tawrat) to Moses, the Psalms (Zābūr) to David, the Gospel to Jesus, and the Quran to Muhammad; those prophets are considered "Messengers" or rasūl. Other main prophets are considered messengers or nabī, even if they didn't receive a Book from God. Examples include the messenger-prophet Aaron| (Hārūn), the messenger-prophet Ishmael (Ismāʿīl)) and the messenger-prophet Joseph (Yūsuf).
|
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Although it offers many incidents from the lives of many prophets, the Quran focuses with special narrative and rhetorical emphasis on the careers of the first four of these five major prophets. Of all the figures before Muhammad, the significance of Jesus in Islam is reflected in his being mentioned in the Quran in 93 verses with various titles attached such as "Son of Mary" and other relational terms, mentioned directly and indirectly, over 187 times.[65][66] He is thus the most mentioned person in the Quran by reference; 25 times by the name Isa, third-person 48 times, first-person 35 times, and the rest as titles and attributes.[67][68][69] Moses(Musa) and Abraham(ibrahim) are also referred to frequently in the Quran. As for the fifth, the Quran is frequently addressed directly to Muhammad, and it often discusses situations encountered by him. Direct use of his name in the text, however, is rare. Rarer still is the mention of Muhammad's contemporaries.
|
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In modern times the term "prophet" can be somewhat controversial. Many Christians with Pentecostal or charismatic beliefs believe in the continuation of the gift of prophecy and the continuation of the role of prophet as taught in Ephesians 4.[70] The content of prophecies can vary widely. Prophecies are often spoken as quotes from God. They may contain quotes from scripture, statements about the past or current situation, or predictions of the future. Prophecies can also 'make manifest the secrets' of the hearts of other people, telling about the details of their lives. Sometimes, more than one person in a congregation will receive the same message in prophecy, with one giving it before another.
|
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|
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Other movements claim to have prophets. In France, Michel Potay says he received a revelation, called The Revelation of Arès, dictated by Jesus in 1974, then by God in 1977. He is considered a prophet by his followers, the Pilgrims of Arès.
|
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The Bahá'í Faith refers to what are commonly called prophets as "Manifestations of God" who are directly linked with the concept of Progressive revelation. Bahá'ís believe that God expresses this will at all times and in many ways, including through a series of divine messengers referred to as "Manifestations of God" or "divine educators".[71] In expressing God's intent, these Manifestations are seen to establish religion in the world. Thus they are seen as an intermediary between God and humanity.[72]
|
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The Manifestations of God are not seen as incarnations of God, and are also not seen as ordinary mortals. Instead, the Bahá'í concept of the Manifestation of God emphasizes simultaneously the humanity of that intermediary and the divinity in the way they show forth the will, knowledge and attributes of God; thus they have both human and divine stations.[72]
|
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|
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+
In addition to the Manifestations of God, there are also minor prophets. While the Manifestations of God, or major prophets, are compared to the Sun (which produces its own heat and light), minor prophets are compared to the Moon (which receives its light from the sun). Moses, for example, is taught as having been a Manifestation of God and his brother Aaron a minor prophet. Moses spoke on behalf of God, and Aaron spoke on behalf of Moses (Exodus 4:14–17).[73] Other Jewish prophets are considered minor prophets, as they are considered to have come in the shadow of the dispensation of Moses to develop and consolidate the process he set in motion.
|
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+
A number of modern catholic saints have been claimed to have powers of prophecy, such as Padre Pio [74] and Alexandrina Maria da Costa.[75]
|
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|
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+
In addition to this many modern Marian apparitions included prophecies in them about the world and about the local areas. The Fátima apparition in 1917 included a prophecy given by Mary to three children, that on October 13, 1917 there would be a great miracle for all to see at Fátima, Portugal, and on that day tens of thousands of people headed to Fátima to see what would happen including newspaper journalists. Many witnesses, including journalists, claimed to see the sun "dance" in the sky in the afternoon of that day, exactly as the visionaries had predicted several months before.[76] The Kibeho apparition in Rwanda in the 1980s included many prophecies about great violence and destruction that was coming, and the Rwandan genocide only ten years later was interpreted by the visionaries as the fulfilment of these prophecies [77]
|
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|
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+
Several miracles and a vision of the identity of the last 112 Popes were attributed to Saint Malachy, the Archbishop of Armagh (1095–1148).
|
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|
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+
Jehovah's Witnesses do not consider any single person in their modern-day organization to be a prophet.[78] Their literature has referred to their organization collectively as God's "prophet" on earth; this is understood, however, in the sense of declaring their interpretation of God's judgments from the Bible along with God's guidance of His Holy Spirit.[79][80] Their publishing company, Watch Tower, and official position magazine, The Watchtower, have asserted: "Ever since The Watchtower began to be published in July 1879 it has looked ahead into the future... No, The Watchtower is no inspired prophet, but it follows and explains a Book of prophecy the predictions in which have proved to be unerring and unfailing till now. The Watchtower is therefore under safe guidance. It may be read with confidence, for its statements may be checked against that prophetic Book."[81] They also claim that they are God's one and only true channel to mankind on earth, and used by God for this purpose.
|
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|
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+
They have made many eschatological forecasts, some of which have led people (including followers) to incorrect assumptions. One example is The Watchtower's assertions that the end of the "Gentile times" or "times of the nations" would occur in 1914; even prominent Watch Tower representatives such as A. H. Macmillan incorrectly concluded and overstated their expectations. As a result, The Watchtower has acknowledged that Jehovah's Witnesses "have made mistakes in their understanding of what would occur at the end of certain time periods."[82] Concurrently with these exceptions, Jehovah's Witnesses in their literature and assemblies have taught their leadership was personally chosen by Jesus Christ in 1919 (a prophetic year in Jehovah's Witnesses eschatology) and that they are "God's sole channel on earth," and "Jehovah's spirit directed organization".
|
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|
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+
Joseph Smith, who established the Church of Christ in 1830, is considered a prophet by members of the Latter Day Saint movement, of which The Church of Jesus Christ of Latter-day Saints (LDS Church) is the largest denomination. Additionally, many churches within the movement believe in a succession of modern prophets (accepted by Latter Day Saints as "prophets, seers, and revelators") since the time of Joseph Smith. Russell M. Nelson is the current Prophet and President of The Church of Jesus Christ of Latter-day Saints.
|
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|
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+
The Seventh-day Adventist Church, which was established in 1863, believes that Ellen G. White, one of the church's founders, was given the spiritual gift of prophecy.
|
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|
78 |
+
The Branch Davidians are a religious cult which was founded in 1959 by Benjamin Roden as an offshoot of the Seventh-Day Adventist Church. David Koresh, who died in the well-known Waco Siege in 1993, claimed to be their final prophet and "the Son of God, the Lamb" in 1983.
|
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+
|
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+
The Ahmadiyya movement in Islam believes that Mirza Ghulam Ahmad was a non law-bearing Prophet, who claimed to be a fulfillment of the various Islamic prophecies regarding the spiritual second advent of Jesus of Nazareth near the end times.[87]
|
81 |
+
|
82 |
+
Nathan of Gaza was a theologian and author who became famous as a prophet for the alleged messiah, Sabbatai Zevi.
|
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+
|
84 |
+
Divination remains an important aspect of the lives of the people of contemporary Africa, especially amongst the usually rural, socially traditionalistic segments of its population. In arguably its most influential manifestation, the system of prophecy practiced by the Babalawos and Iyanifas of the historically Yoruba regions of West Africa have bequeathed to the world a corpus of fortune-telling poetic methodologies so intricate that they have been added by UNESCO to its official intangible cultural heritage of the World list.
|
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+
|
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+
Tenrikyo's prophet, Nakayama Miki, is believed by Tenrikyoans to have been a messenger of God.[88]
|
87 |
+
|
88 |
+
The Great Peacemaker (sometimes referred to as Deganawida or Dekanawida) co-founded the Haudenosaunee league in pre-Columbian times. In retrospect, his prophecy of the boy seer could appear to refer to the conflict between natives and Europeans (white serpent).
|
89 |
+
|
90 |
+
From 1805 until the Battle of Tippecanoe that falsified his predictions in 1811, the "Shawnee prophet" Tenskwatawa lead an Indian alliance to stop Europeans to take more and more land going west. He reported visions he had. He is said to have accurately predicted a solar eclipse. His brother Tecumseh re-established the alliance for Tecumseh's War, that ended with the latter's death in 1813. Tecumseh fought together with British forces that, in the area of the Great Lakes, occupied essentially today's territory of Canada.
|
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|
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+
Francis the Prophet, influenced by Tecumseh and Tenskwatawa, was a leader of the Red Stick faction of the Creek Indians. He traveled to England in 1815 as a representative of the "four Indian nations" in an unsuccessful attempt to get Great Britain to help them resist the expansionism of the white settlers.
|
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+
|
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+
20 years later (1832), Wabokieshiek, the "Winnebago Prophet", after whom Prophetstown has been named, (also called "White Cloud") claimed that British forces would support the Indians in the Black Hawk War against the United States as 20 years earlier (based on "visions"). They did not, and no longer he was considered a "prophet".
|
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+
|
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+
In 1869, the Paiute Wodziwob founded the Ghost Dance movement. The dance rituals were an occasion to announce his visions of an earthquake that would swallow the whites. He seems to have died in 1872.
|
97 |
+
|
98 |
+
The Northern Paiute Wovoka claimed he had a vision during the solar eclipse of January 1, 1889, that the Paiute dead would come back and the whites would vanish from America, provided the natives performed Ghost Dances. This idea spread among other Native American peoples. The government were worried about a rebellion and sent troops, which lead to the death of Sitting Bull and to the Wounded Knee massacre in 1890.
|
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+
|
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+
Clifford Trafzer compiled an anthology of essays on the topic, American Indian Prophets. Trafzer, Clifford. American Indian Prophets: Religious Leaders and Revitalization Movements. Sierra Oaks Publishing Co. ISBN 9780940113022.
|
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|
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+
In the late 20th century the appellation of prophet has been used to refer to individuals particularly successful at analysis in the field of economics, such as in the derogatory prophet of greed. Alternatively, social commentators who suggest escalating crisis are often called prophets of doom.[90][91]
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Prose is a form or technique of language that exhibits a natural flow of speech and grammatical structure. Novels, textbooks and newspaper articles are all examples of prose. The word prose is frequently used in opposition to traditional poetry, which is language with a regular structure and a common unit of verse based on metre or rhyme. However, as T. S. Eliot noted, whereas "the distinction between verse and prose is clear, the distinction between poetry and prose is obscure";[1] developments in modern literature, including free verse and prose poetry, have led to the two techniques indicating two ends on a spectrum of ways to compose language, as opposed to two discrete options.
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Prose in its simplicity and loosely defined structure is broadly adaptable to spoken dialogue, factual discourse, and to topical and fictional writing. It is systematically produced and published within literature, journalism (including newspapers, magazines, and broadcasting), encyclopedias, film, history, philosophy, law, and in almost all forms and processes requiring human communication.
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The word "prose" first appears in English in the 14th century. It is derived from the Old French prose, which in turn originates in the Latin expression prosa oratio (literally, straightforward or direct speech).[2]
|
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|
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+
Isaac Newton in The Chronology of Ancient Kingdoms wrote "The Greek Antiquities are full of Poetical Fictions, because the Greeks wrote nothing in Prose, before the Conquest of Asia by Cyrus the Persian. Then Pherecydes Scyrius and Cadmus Milesius introduced the writing in Prose."[3]
|
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|
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Prose lacks the more formal metrical structure of verse that can be found in traditional poetry. Prose comprises full grammatical sentences, which then constitute paragraphs while overlooking aesthetic appeal, whereas poetry often involves a metrical or rhyming scheme. Some works of prose contain traces of metrical structure or versification and a conscious blend of the two literature formats known as prose poetry. Verse is considered to be more systematic or formulaic, whereas prose is the most reflective of ordinary (often conversational) speech. On this point, Samuel Taylor Coleridge jokingly requested that novice poets should know the "definitions of prose and poetry; that is, prose—words in their best order; poetry—the best words in their best order."[4]
|
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Monsieur Jourdain asked for something to be written in neither verse nor prose. A philosophy master replied there is no other way to express oneself than with prose or verse, for the simple reason being that everything that is not prose is verse, and everything that is not verse is prose. Molière, Le Bourgeois gentilhomme[5]
|
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|
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+
... I believe a story can be wrecked by a faulty rhythm in a sentence— especially if it occurs toward the end—or a mistake in paragraphing, even punctuation. Henry James is the maestro of the semicolon. Hemingway is a first-rate paragrapher. From the point of view of ear, Virginia Woolf never wrote a bad sentence. I don't mean to imply that I successfully practice what I preach. I try, that's all.[6]
|
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+
|
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+
Many types of prose exist, which include nonfictional prose, heroic prose,[7] prose poem,[8] polyphonic prose, alliterative prose, prose fiction, and village prose in Russian literature.[9] A prose poem is a composition in prose that has some of the qualities of a poem.[10]
|
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|
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Many forms of creative or literary writing use prose, including novels and short stories. Writer Truman Capote thought that the short story was "the most difficult and disciplining form of prose writing extant".[6]
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The prostate is an exocrine gland of the male reproductive system in most mammals and some invertebrates. It differs considerably among species anatomically, chemically, and physiologically. The word prostate comes from Ancient Greek προστάτης, prostátēs, meaning "one who stands before", "protector", "guardian".[1]
|
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|
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+
Anatomically, the prostate is found below the bladder, with the urethra passing through it. It is described in gross anatomy as consisting of lobes, and in microanatomy by zone. It does not have a capsule; rather an integral fibromuscular band surrounds it.[2] It is sheathed in the muscles of the pelvic floor, which contract during ejaculation. The prostate contains some smooth muscles that help expel semen during ejaculation.
|
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|
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+
The function of the prostate is to secrete a fluid which contributes to the volume of the semen. This prostatic fluid is slightly alkaline, milky or white in appearance, and in humans usually constitutes roughly 30% of the volume of semen, the other 70% being spermatozoa and seminal vesicle fluid.[3] The alkalinity of semen helps neutralize the acidity of the vaginal tract, prolonging the lifespan of sperm.
|
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|
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+
The prostatic fluid is expelled in the first part of ejaculate, together with most of the sperm. In comparison with the few spermatozoa expelled together with mainly seminal vesicular fluid, those in prostatic fluid have better motility, longer survival, and better protection of genetic material.
|
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|
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+
Disorders of the prostate include enlargement, inflammation, infection, and cancer.
|
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|
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+
The prostate is a gland of the male reproductive system. In adults, it is about the size of a walnut,[4] and has an average weight of about 11 grams, usually ranging between 7 and 16 grams.[5] The prostate is located in the pelvis. It sits below the urinary bladder and surrounds the urethra. The part of the urethra passing through it is called the prostatic urethra, which joints with the two ejaculatory ducts.[4] The prostate is covered in a surface called the prostatic capsule or prostatic fascia.[6]
|
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|
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+
The internal structure of the prostate has been described using both lobes and zones.[7][4] Because of the variation in descriptions and definitions of lobes, the zone classification is used more predominantly.[4]
|
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+
|
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+
The prostate has been described as consisting of three or four zones.[4][6] Zones are more typically able to be seen on histology, or in medical imaging, such as ultrasound or MRI.[4][7] The zones are:
|
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+
|
19 |
+
The "lobe" classification describes lobes that, whilst originally defined in the foetus, are also visible in gross anatomy, including dissection and when viewed endoscopically.[7][6] The five lobes are the anterior lobe or isthmus; the posterior lobe; the right and left lateral lobes, and the middle or median lobe.
|
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|
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+
Lobes of prostate
|
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|
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+
Zones of prostate
|
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|
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+
The prostate receives blood through the inferior vesical artery, internal pudendal artery, and middle rectal arteries. These vessels enter the prostate on its outer posterior surface where it meets the bladder, and travel forward to the apex of the prostate.[6] Both the inferior vesicle and the middle rectal arteries often arise together directly from the internal iliac arteries. On entering the bladder, the inferior vesical artery splits into a urethral branch, supplying the urethral prostate; and a capsular branch, which travels around the capsule and has smaller branches which perforate into the prostate.[6]
|
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|
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+
The veins of the prostate form a network – the prostatic venous plexus, primarily around its front and outer surface.[6] This network also receives blood from the deep dorsal vein of the penis, and is connected via branches to the vesical plexus and internal pudendal veins.[6] Veins drain into the vesical and then internal iliac veins.[6]
|
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|
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+
The lymphatic drainage of the prostate depends on the positioning of the area. Vessels surrounding the vas deferens, some of the vessels in the seminal vesicle, and a vessel from the posterior surface of the prostate drain into the external iliac lymph nodes.[6] Some of the seminal vesicle vessels, prostatic vessels, and vessels from the anterior prostate drain into internal iliac lymph nodes.[6] Vessels of the prostate itself also drain into the obturator and sacral lymph nodes.[6]
|
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|
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+
Imaging showing the inferior vesical, inferior pudendal and middle rectal arteries arising from the internal iliac arteries.
|
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+
|
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+
Image showing the external iliac lymph nodes and their positions around the external iliac artery and vein
|
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+
|
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+
The prostate consists of glandular and connective tissue.[4] Tall column-shaped cells form the lining (the epithelium) of the glands.[4] These form one layer or may be pseudostratified.[6] The epithelium is highly variable and areas of low cuboidal or flat cells can also be present, with transitional epithelium in the outer regions of the longer ducts.[11] The glands are formed as many follicles, which in drain into canals and subsequently 12–20 main ducts, These in turn drain into the urethra as it passes through the prostate.[6] There are also a small amount of flat cells, which sit next to the basement membranes of glands, and act as stem cells.[4]
|
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|
37 |
+
The connective tissue of the prostate is made up of fibrous tissue and smooth muscle.[4] The fibrous tissue separates the gland into lobules.[4] It also sits between the glands and is composed of randomly orientated smooth-muscle bundles that are continuous with the bladder.[12]
|
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|
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+
Over time, thickened secretions called corpora amylacea accumulate in the gland.[4]
|
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+
|
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+
Microscopic glands of the prostate
|
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+
|
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+
About 20,000 protein coding genes are expressed in human cells and almost 75% of these genes are expressed in the normal prostate.[13][14] About 150 of these genes are more specifically expressed in the prostate with about 20 genes being highly prostate specific.[15] The corresponding specific proteins are expressed in the glandular and secretory cells of the prostatic gland and have functions that are important for the characteristics of semen. Some of the prostate specific proteins are enzymes, such as the prostate specific antigen (PSA), and the ACPP protein.
|
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+
|
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+
In the developing embryo, at the hind end lies an inpouching called the cloaca. This, over the fourth to the seventh week, divides into a urogenital sinus and the beginnings of the anal canal, with a wall forming between these two inpouchings called the urorectal septum.[16] The urogenital sinus divides into three parts, with the middle part forming the urethra; the upper part is largest and becomes the urinary bladder, and the lower part then changes depending on the biological sex of the embryo.[16]
|
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+
|
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+
The prostatic part of the urethra develops from the middle, pelvic, part of the urogenital sinus, which is of endodermal origin.[17] Around the end of the third month of embryonic life, outgrowths arise from the prostatic part of the urethra and grow into the surrounding mesenchyme.[17] The cells lining this part of the urethra differentiate into the glandular epithelium of the prostate.[17] The associated mesenchyme differentiates into the dense connective tissue and the smooth muscle of the prostate.[18]
|
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+
|
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+
Condensation of mesenchyme, urethra, and Wolffian ducts gives rise to the adult prostate gland, a composite organ made up of several tightly fused glandular and non-glandular components. To function properly, the prostate needs male hormones (androgens), which are responsible for male sex characteristics. The main male hormone is testosterone, which is produced mainly by the testicles. It is dihydrotestosterone (DHT), a metabolite of testosterone, that predominantly regulates the prostate. The prostate gland enlarges over time, until the fourth decade of life.[6]
|
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|
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+
The prostate secretes fluid which becomes part of semen. Semen is the fluid emitted (ejaculated) by males during the sexual response.[19] When sperm is emitted, it is transmitted from the vas deferens into the male urethra via the ejaculatory ducts, which lie within the prostate gland.[19] Ejaculation is the expulsion of semen from the urethra.[19] Semen is moved into the urethra following contractions of the smooth muscle of the vas deferens and seminal vesicles, following stimulation, primarily of the glans penis. Stimulation sends nerve signals via the internal pudendal nerves to the upper lumbar spine; the nerve signals causing contraction act via the hypogastric nerves.[19] After traveling into the urethra, the seminal fluid is ejaculated by contraction of the bulbocavernosus muscle.[19] The secretions of the prostate include proteolytic enzymes, prostatic acid phosphatase, fibrinolysin, zinc, and prostate-specific antigen.[6] Together with the secretions from the seminal vesicles, these form the major fluid part of semen.[6]
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It is possible for some men to achieve orgasm solely through stimulation of the prostate gland, such as prostate massage or anal intercourse.[20][21]
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Prostatitis is inflammation of the prostate gland. It can be caused by infection with bacteria, or other noninfective causes. Inflammation of the prostate can cause painful urination or ejaculation, groin pain, difficulty passing urine, or constitutional symptoms such as fever or tiredness.[22] When inflamed, the prostate becomes enlarged and is tender when touched during digital rectal examination. A culprit bacteria may grow in a urine culture.[22]
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Acute prostatitis and chronic bacterial prostatitis are treated with antibiotics.[22] Chronic non-bacterial prostatitis, or male chronic pelvic pain syndrome is treated by a large variety of modalities including the medications alpha blockers, nonsteroidal antiinflammatories and amitriptyline,[22] antihistamines, and other anxiolytics.[23] Other treatments that are not medications may include physical therapy,[24] psychotherapy, nerve modulators, and surgery. More recently, a combination of trigger point and psychological therapy has proved effective for category III prostatitis as well.[23]
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An enlarged prostate is called prostatomegaly, with benign prostatic hyperplasia (BPH) being the most common cause. BPH refers to an enlargement of the prostate due to an increase in the number of cells that make up the prostate (hyperplasia) from a cause that is not a malignancy. It is very common in older men.[22] It is often diagnosed when the prostate has enlarged to the point where urination becomes difficult. Symptoms include needing to urinate often (urinary frequency) or taking a while to get started (urinary hesitancy). If the prostate grows too large, it may constrict the urethra and impede the flow of urine, making urination painful and difficult, or in extreme cases completely impossible, causing urinary retention.[22] Over time, chronic retention may cause the bladder to become larger and cause a backflow of urine into the kidneys (hydronephrosis).[22]
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BPH can be treated with medication, a minimally invasive procedure or, in extreme cases, surgery that removes the prostate. In general, treatment often begins with an alpha-1 adrenergic receptor antagonist medication such as tamsulosin, which reduces the tone of the smooth muscle found in the urethra that passes through the prostate, making it easier for urine to pass through.[22] For people with persistent symptoms, procedures may be considered. The surgery most often used in such cases is transurethral resection of the prostate,[22] in which an instrument is inserted through the urethra to remove prostate tissue that is pressing against the upper part of the urethra and restricting the flow of urine. Minimally invasive procedures include transurethral needle ablation of the prostate and transurethral microwave thermotherapy.[25] These outpatient procedures may be followed by the insertion of a temporary stent, to allow normal voluntary urination, without exacerbating irritative symptoms.[26]
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Prostate cancer is one of the most common cancers affecting older men in the UK, US, Northern Europe and Australia, and a significant cause of death for elderly men worldwide.[27] Often, a person does not have symptoms; when they do occur, symptoms may include urinary frequency, urgency, hesitation and other symptoms associated with BPH. Uncommonly, such cancers may cause weight loss, retention of urine, or symptoms such as back pain due to metastatic lesions that have spread outside of the prostate.[22]
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A digital rectal examination and the measurement of a prostate specific antigen (PSA) level are usually the first investigations done to check for prostate cancer. PSA values are difficult to interpret, because a high value might be present in a person without cancer, and a low value can be present in someone with cancer.[22] The next form of testing is often the taking of a biopsy to assess for tumour activity and invasiveness.[22] Because of the significant risk of overdiagnosis with widespread screening in the general population, prostate cancer screening is controversial.[28] If a tumour is confirmed, medical imaging such as an MRI or bone scan may be done to check for the presence of tumour metastases in other parts of the body.[22]
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Prostate cancer that is only present in the prostate is often treated with either surgical removal of the prostate or with radiotherapy or by the insertion of small radioactive particles (brachytherapy)[22] Cancer that has spread to other parts of the body is usually treated also with hormone therapy, to deprive a tumour of sex hormones (androgens) that stimulate proliferation. This is often done through the use of GnRH analogues or agents that block the receptors that androgens act at, such as bicalutamide; occasionally, surgical removal of the testes may be done instead.[22] Cancer that does not respond to hormonal treatment, or that progresses after treatment, might be treated with chemotherapy such as docetaxel. Radiotherapy may also be used to help with pain associated with bony lesions.[22]
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Sometimes, the decision may be made not to treat prostate cancer. If a cancer is small and localised, the decision may be made to monitor for cancer activity at intervals ("active surveillance") and defer treatment.[22] If a person, because of frailty or other medical conditions or reasons, has a life expectancy less than ten years, then the impacts of treatment may outweigh any perceived benefits.[22]
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The prostate was first formally identified by Venetian anatomist Niccolò Massa in Anatomiae libri introductorius (Introduction to Anatomy) 1536 and illustrated by Flemish anatomist Andreas Vesalius in Tabulae anatomicae sex (six anatomical tables) in 1538.[29][7] Massa described it as a "glandular flesh upon which rests the neck of the bladder," and Vesalius as a "glandular body".[30] The first time a word similar to 'prostate' was used to describe the gland is credited to André du Laurens in 1600, who described it as a term used by anatomists at the time.[30][7] The term was however used at least as early as 1549 by French surgeon Ambroise Pare.[7]
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At the time, Du Laurens was describing what was considered to be a pair of organs (not the single two-lobed organ), and the Latin term prostatae that was used was a mistranslation of the term for the Ancient Greek word used to describe the seminal vesicles, parastatai;[30] although it has been argued that surgeons in Ancient Greece and Rome must have at least seen the prostate as an anatomical entity.[7] The term prostatae was taken rather than the grammatically correct prostator (singular) and prostatores (plural) because the gender of the Ancient Greek term was taken as female, when it was in fact male.[30]
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The fact that the prostate was one and not two organs was an idea popularised throughout the early 18th century, as was the English language term used to describe the organ, prostate.[30] A monograph, "Practical observations on the treatment of the diseases of the prostate gland" by Everard Home in 1811 was important in the history of the prostate by describing and naming anatomical parts of the prostate.[30] The idea of the lobes of the prostate was popularised following anatomical studies conducted by American urologist Oswald Lowsley in 1912.[7] John E. McNeal first proposed the idea of "zones" in 1968. McNeal found that the relatively homogeneous cut surface of an adult prostate in no way resembled "lobes" and thus led to the description of "zones".[31]
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Prostate cancer was first described in a speech to the Medical and Chiurgical Society of London in 1853 by surgeon John Adams[32][33] and increasingly described by the late 19th century.[34] Prostate cancer was initially considered a rare disease, probably because of shorter life expectancies and poorer detection methods in the 19th century. The first treatments of prostate cancer were surgeries to relieve urinary obstruction.[35] Removal of the gland was first described in 1851, by Samuel David Gross,[36][34] and removal for prostate cancer (radical perineal prostatectomy) was first performed in 1904 by Hugh H. Young at Johns Hopkins Hospital.[37][34] Surgical removal of the testes (orchiectomy) to treat prostate cancer was first performed in the 1890s, but with limited success. Transurethral resection of the prostate (TURP) replaced radical prostatectomy for symptomatic relief of obstruction in the middle of the 20th century because it could better preserve penile erectile function. Radical retropubic prostatectomy was developed in 1983 by Patrick Walsh.[38] This surgical approach allowed for removal of the prostate and lymph nodes with maintenance of penile function. In 1941, Charles B. Huggins published studies in which he used estrogen to oppose testosterone production in men with metastatic prostate cancer. This discovery of "chemical castration" won Huggins the 1966 Nobel Prize in Physiology or Medicine.[39]
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The role of the gonadotropin-releasing hormone (GnRH) in reproduction was determined by Andrzej W. Schally and Roger Guillemin, who both won the 1977 Nobel Prize in Physiology or Medicine for this work. GnRH receptor agonists, such as leuprorelin and goserelin, were subsequently developed and used to treat prostate cancer.[40][41] Radiation therapy for prostate cancer was first developed in the early 20th century and initially consisted of intraprostatic radium implants. External beam radiotherapy became more popular as stronger X-ray radiation sources became available in the middle of the 20th century. Brachytherapy with implanted seeds (for prostate cancer) was first described in 1983.[42] Systemic chemotherapy for prostate cancer was first studied in the 1970s. The initial regimen of cyclophosphamide and 5-fluorouracil was quickly joined by multiple regimens using a host of other systemic chemotherapy drugs.[43]
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The prostate is found as a male accessory gland in all placental mammals excepting xenarthrans, martens, badgers, and otters.[44] The prostate glands of male marsupials are disseminate[45] and proportionally larger than those of placental mammals.[46] In some marsupial species, the size of the prostate gland changes seasonally.[47] The structure of the prostate varies, ranging from tubuloalveolar (as in humans) to branched tubular. The gland is particularly well developed in dogs, foxes and boars, though in other mammals, such as bulls, it can be small and inconspicuous.[48][49] The prostate is the only accessory gland that occurs in male dogs.[50] Dogs can produce in one hour as much prostatic fluid as a human can in a day. They excrete this fluid along with their urine to mark their territory.[51] In many rodents and bats, the prostatic fluid contains a coagulant. This mixes with and coagulates semen during copulation to form a mating plug that temporarily prevents further copulation.[52][53] In cetaceans the prostate is composed of diffuse urethral glands[54] and is surrounded by a very powerful compressor muscle.[55]
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Prostatic secretions vary among species. They are generally composed of simple sugars and are often slightly alkaline.[56]
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The prostate gland originates with tissues in the urethral wall. This means the urethra, a compressible tube used for urination, runs through the middle of the prostate. This leads to an evolutionary design fault for some mammals, including human males. The prostate is prone to infection and enlargement later in life, constricting the urethra so urinating becomes slow and painful.[57]
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Monotremes and marsupial moles lack prostates, instead having simpler cloacal glands that carry out their function.[58][59]
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A prostate gland occurs in some invertebrate species, such as gastropods (slugs and snails).[60]
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The Skene's gland, also known as the female prostate, is found in both female humans and rodents. Historically it was thought to be a vestigial organ, but it has been discovered that it produces the same protein markers, PSA and PAB, as the male prostate.[61] It functions as a histologic homolog to the male prostate gland.[62][63]
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Media related to Prostate at Wikimedia Commons
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en/4828.html.txt
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1 |
+
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Proteins are large biomolecules, or macromolecules, consisting of one or more long chains of amino acid residues. Proteins perform a vast array of functions within organisms, including catalysing metabolic reactions, DNA replication, responding to stimuli, providing structure to cells, and organisms, and transporting molecules from one location to another. Proteins differ from one another primarily in their sequence of amino acids, which is dictated by the nucleotide sequence of their genes, and which usually results in protein folding into a specific 3D structure that determines its activity.
|
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A linear chain of amino acid residues is called a polypeptide. A protein contains at least one long polypeptide. Short polypeptides, containing less than 20–30 residues, are rarely considered to be proteins and are commonly called peptides, or sometimes oligopeptides. The individual amino acid residues are bonded together by peptide bonds and adjacent amino acid residues. The sequence of amino acid residues in a protein is defined by the sequence of a gene, which is encoded in the genetic code. In general, the genetic code specifies 20 standard amino acids; but in certain organisms the genetic code can include selenocysteine and—in certain archaea—pyrrolysine. Shortly after or even during synthesis, the residues in a protein are often chemically modified by post-translational modification, which alters the physical and chemical properties, folding, stability, activity, and ultimately, the function of the proteins. Some proteins have non-peptide groups attached, which can be called prosthetic groups or cofactors. Proteins can also work together to achieve a particular function, and they often associate to form stable protein complexes.
|
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Once formed, proteins only exist for a certain period and are then degraded and recycled by the cell's machinery through the process of protein turnover. A protein's lifespan is measured in terms of its half-life and covers a wide range. They can exist for minutes or years with an average lifespan of 1–2 days in mammalian cells. Abnormal or misfolded proteins are degraded more rapidly either due to being targeted for destruction or due to being unstable.
|
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|
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Like other biological macromolecules such as polysaccharides and nucleic acids, proteins are essential parts of organisms and participate in virtually every process within cells. Many proteins are enzymes that catalyse biochemical reactions and are vital to metabolism. Proteins also have structural or mechanical functions, such as actin and myosin in muscle and the proteins in the cytoskeleton, which form a system of scaffolding that maintains cell shape. Other proteins are important in cell signaling, immune responses, cell adhesion, and the cell cycle. In animals, proteins are needed in the diet to provide the essential amino acids that cannot be synthesized. Digestion breaks the proteins down for use in the metabolism.
|
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|
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Proteins may be purified from other cellular components using a variety of techniques such as ultracentrifugation, precipitation, electrophoresis, and chromatography; the advent of genetic engineering has made possible a number of methods to facilitate purification. Methods commonly used to study protein structure and function include immunohistochemistry, site-directed mutagenesis, X-ray crystallography, nuclear magnetic resonance and mass spectrometry.
|
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Proteins were recognized as a distinct class of biological molecules in the eighteenth century by Antoine Fourcroy and others, distinguished by the molecules' ability to coagulate or flocculate under treatments with heat or acid.[1] Noted examples at the time included albumin from egg whites, blood serum albumin, fibrin, and wheat gluten.
|
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|
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+
Proteins were first described by the Dutch chemist Gerardus Johannes Mulder and named by the Swedish chemist Jöns Jacob Berzelius in 1838.[2][3] Mulder carried out elemental analysis of common proteins and found that nearly all proteins had the same empirical formula, C400H620N100O120P1S1.[4] He came to the erroneous conclusion that they might be composed of a single type of (very large) molecule. The term "protein" to describe these molecules was proposed by Mulder's associate Berzelius; protein is derived from the Greek word πρώτειος (proteios), meaning "primary",[5] "in the lead", or "standing in front",[6] + -in. Mulder went on to identify the products of protein degradation such as the amino acid leucine for which he found a (nearly correct) molecular weight of 131 Da.[4] Prior to "protein", other names were used, like "albumins" or "albuminous materials" (Eiweisskörper, in German).[7]
|
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|
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+
Early nutritional scientists such as the German Carl von Voit believed that protein was the most important nutrient for maintaining the structure of the body, because it was generally believed that "flesh makes flesh."[8] Karl Heinrich Ritthausen extended known protein forms with the identification of glutamic acid. At the Connecticut Agricultural Experiment Station a detailed review of the vegetable proteins was compiled by Thomas Burr Osborne. Working with Lafayette Mendel and applying Liebig's law of the minimum in feeding laboratory rats, the nutritionally essential amino acids were established. The work was continued and communicated by William Cumming Rose. The understanding of proteins as polypeptides came through the work of Franz Hofmeister and Hermann Emil Fischer in 1902.[9][10] The central role of proteins as enzymes in living organisms was not fully appreciated until 1926, when James B. Sumner showed that the enzyme urease was in fact a protein.[11]
|
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|
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+
The difficulty in purifying proteins in large quantities made them very difficult for early protein biochemists to study. Hence, early studies focused on proteins that could be purified in large quantities, e.g., those of blood, egg white, various toxins, and digestive/metabolic enzymes obtained from slaughterhouses. In the 1950s, the Armour Hot Dog Co. purified 1 kg of pure bovine pancreatic ribonuclease A and made it freely available to scientists; this gesture helped ribonuclease A become a major target for biochemical study for the following decades.[4]
|
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Linus Pauling is credited with the successful prediction of regular protein secondary structures based on hydrogen bonding, an idea first put forth by William Astbury in 1933.[12] Later work by Walter Kauzmann on denaturation,[13][14] based partly on previous studies by Kaj Linderstrøm-Lang,[15] contributed an understanding of protein folding and structure mediated by hydrophobic interactions.
|
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+
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+
The first protein to be sequenced was insulin, by Frederick Sanger, in 1949. Sanger correctly determined the amino acid sequence of insulin, thus conclusively demonstrating that proteins consisted of linear polymers of amino acids rather than branched chains, colloids, or cyclols.[16] He won the Nobel Prize for this achievement in 1958.[17]
|
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+
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+
The first protein structures to be solved were hemoglobin and myoglobin, by Max Perutz and Sir John Cowdery Kendrew, respectively, in 1958.[18][19] As of 2017[update], the Protein Data Bank has over 126,060 atomic-resolution structures of proteins.[20] In more recent times, cryo-electron microscopy of large macromolecular assemblies[21] and computational protein structure prediction of small protein domains[22] are two methods approaching atomic resolution.
|
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+
Most proteins consist of linear polymers built from series of up to 20 different L-α- amino acids. All proteinogenic amino acids possess common structural features, including an α-carbon to which an amino group, a carboxyl group, and a variable side chain are bonded. Only proline differs from this basic structure as it contains an unusual ring to the N-end amine group, which forces the CO–NH amide moiety into a fixed conformation.[23] The side chains of the standard amino acids, detailed in the list of standard amino acids, have a great variety of chemical structures and properties; it is the combined effect of all of the amino acid side chains in a protein that ultimately determines its three-dimensional structure and its chemical reactivity.[24]
|
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+
The amino acids in a polypeptide chain are linked by peptide bonds. Once linked in the protein chain, an individual amino acid is called a residue, and the linked series of carbon, nitrogen, and oxygen atoms are known as the main chain or protein backbone.[25]
|
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+
|
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+
The peptide bond has two resonance forms that contribute some double-bond character and inhibit rotation around its axis, so that the alpha carbons are roughly coplanar. The other two dihedral angles in the peptide bond determine the local shape assumed by the protein backbone.[26] The end with a free amino group is known as the N-terminus or amino terminus, whereas the end of the protein with a free carboxyl group is known as the C-terminus or carboxy terminus (the sequence of the protein is written from N-terminus to C-terminus, from left to right).
|
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+
|
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+
The words protein, polypeptide, and peptide are a little ambiguous and can overlap in meaning. Protein is generally used to refer to the complete biological molecule in a stable conformation, whereas peptide is generally reserved for a short amino acid oligomers often lacking a stable 3D structure. But the boundary between the two is not well defined and usually lies near 20–30 residues.[27] Polypeptide can refer to any single linear chain of amino acids, usually regardless of length, but often implies an absence of a defined conformation.
|
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|
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+
Proteins can interact with many types of molecules, including with other proteins, with lipids, with carboyhydrates, and with DNA.[28][29][30][31]
|
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|
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+
It has been estimated that average-sized bacteria contain about 2 million proteins per cell (e.g. E. coli and Staphylococcus aureus). Smaller bacteria, such as Mycoplasma or spirochetes contain fewer molecules, on the order of 50,000 to 1 million. By contrast, eukaryotic cells are larger and thus contain much more protein. For instance, yeast cells have been estimated to contain about 50 million proteins and human cells on the order of 1 to 3 billion.[32] The concentration of individual protein copies ranges from a few molecules per cell up to 20 million.[33] Not all genes coding proteins are expressed in most cells and their number depends on, for example, cell type and external stimuli. For instance, of the 20,000 or so proteins encoded by the human genome, only 6,000 are detected in lymphoblastoid cells.[34] Moreover, the number of proteins the genome encodes correlates well with the organism complexity. Eukaryotes have 15,000, bacteria have 3,200, archaea have 2,400, and viruses have 42 proteins on average coded in their respective genomes.[35]
|
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+
Proteins are assembled from amino acids using information encoded in genes. Each protein has its own unique amino acid sequence that is specified by the nucleotide sequence of the gene encoding this protein. The genetic code is a set of three-nucleotide sets called codons and each three-nucleotide combination designates an amino acid, for example AUG (adenine–uracil–guanine) is the code for methionine. Because DNA contains four nucleotides, the total number of possible codons is 64; hence, there is some redundancy in the genetic code, with some amino acids specified by more than one codon.[36] Genes encoded in DNA are first transcribed into pre-messenger RNA (mRNA) by proteins such as RNA polymerase. Most organisms then process the pre-mRNA (also known as a primary transcript) using various forms of Post-transcriptional modification to form the mature mRNA, which is then used as a template for protein synthesis by the ribosome. In prokaryotes the mRNA may either be used as soon as it is produced, or be bound by a ribosome after having moved away from the nucleoid. In contrast, eukaryotes make mRNA in the cell nucleus and then translocate it across the nuclear membrane into the cytoplasm, where protein synthesis then takes place. The rate of protein synthesis is higher in prokaryotes than eukaryotes and can reach up to 20 amino acids per second.[37]
|
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|
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+
The process of synthesizing a protein from an mRNA template is known as translation. The mRNA is loaded onto the ribosome and is read three nucleotides at a time by matching each codon to its base pairing anticodon located on a transfer RNA molecule, which carries the amino acid corresponding to the codon it recognizes. The enzyme aminoacyl tRNA synthetase "charges" the tRNA molecules with the correct amino acids. The growing polypeptide is often termed the nascent chain. Proteins are always biosynthesized from N-terminus to C-terminus.[36]
|
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+
|
42 |
+
The size of a synthesized protein can be measured by the number of amino acids it contains and by its total molecular mass, which is normally reported in units of daltons (synonymous with atomic mass units), or the derivative unit kilodalton (kDa). The average size of a protein increases from Archaea to Bacteria to Eukaryote (283, 311, 438 residues and 31, 34, 49 kDa respectively) due to a bigger number of protein domains constituting proteins in higher organisms.[35] For instance, yeast proteins are on average 466 amino acids long and 53 kDa in mass.[27] The largest known proteins are the titins, a component of the muscle sarcomere, with a molecular mass of almost 3,000 kDa and a total length of almost 27,000 amino acids.[38]
|
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+
|
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+
Short proteins can also be synthesized chemically by a family of methods known as peptide synthesis, which rely on organic synthesis techniques such as chemical ligation to produce peptides in high yield.[39] Chemical synthesis allows for the introduction of non-natural amino acids into polypeptide chains, such as attachment of fluorescent probes to amino acid side chains.[40] These methods are useful in laboratory biochemistry and cell biology, though generally not for commercial applications. Chemical synthesis is inefficient for polypeptides longer than about 300 amino acids, and the synthesized proteins may not readily assume their native tertiary structure. Most chemical synthesis methods proceed from C-terminus to N-terminus, opposite the biological reaction.[41]
|
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+
|
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+
Most proteins fold into unique 3D structures. The shape into which a protein naturally folds is known as its native conformation.[42] Although many proteins can fold unassisted, simply through the chemical properties of their amino acids, others require the aid of molecular chaperones to fold into their native states.[43] Biochemists often refer to four distinct aspects of a protein's structure:[44]
|
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+
|
48 |
+
Proteins are not entirely rigid molecules. In addition to these levels of structure, proteins may shift between several related structures while they perform their functions. In the context of these functional rearrangements, these tertiary or quaternary structures are usually referred to as "conformations", and transitions between them are called conformational changes. Such changes are often induced by the binding of a substrate molecule to an enzyme's active site, or the physical region of the protein that participates in chemical catalysis. In solution proteins also undergo variation in structure through thermal vibration and the collision with other molecules.[45]
|
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+
|
50 |
+
Proteins can be informally divided into three main classes, which correlate with typical tertiary structures: globular proteins, fibrous proteins, and membrane proteins. Almost all globular proteins are soluble and many are enzymes. Fibrous proteins are often structural, such as collagen, the major component of connective tissue, or keratin, the protein component of hair and nails. Membrane proteins often serve as receptors or provide channels for polar or charged molecules to pass through the cell membrane.[46]
|
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+
A special case of intramolecular hydrogen bonds within proteins, poorly shielded from water attack and hence promoting their own dehydration, are called dehydrons.[47]
|
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+
|
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+
Many proteins are composed of several protein domains, i.e. segments of a protein that fold into distinct structural units. Domains usually also have specific functions, such as enzymatic activities (e.g. kinase) or they serve as binding modules (e.g. the SH3 domain binds to proline-rich sequences in other proteins).
|
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+
|
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+
Short amino acid sequences within proteins often act as recognition sites for other proteins.[48] For instance, SH3 domains typically bind to short PxxP motifs (i.e. 2 prolines [P], separated by two unspecified amino acids [x], although the surrounding amino acids may determine the exact binding specificity). Many such motifs has been collected in the Eukaryotic Linear Motif (ELM) database.
|
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+
|
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+
Proteins are the chief actors within the cell, said to be carrying out the duties specified by the information encoded in genes.[27] With the exception of certain types of RNA, most other biological molecules are relatively inert elements upon which proteins act. Proteins make up half the dry weight of an Escherichia coli cell, whereas other macromolecules such as DNA and RNA make up only 3% and 20%, respectively.[49] The set of proteins expressed in a particular cell or cell type is known as its proteome.
|
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The chief characteristic of proteins that also allows their diverse set of functions is their ability to bind other molecules specifically and tightly. The region of the protein responsible for binding another molecule is known as the binding site and is often a depression or "pocket" on the molecular surface. This binding ability is mediated by the tertiary structure of the protein, which defines the binding site pocket, and by the chemical properties of the surrounding amino acids' side chains. Protein binding can be extraordinarily tight and specific; for example, the ribonuclease inhibitor protein binds to human angiogenin with a sub-femtomolar dissociation constant (<10−15 M) but does not bind at all to its amphibian homolog onconase (>1 M). Extremely minor chemical changes such as the addition of a single methyl group to a binding partner can sometimes suffice to nearly eliminate binding; for example, the aminoacyl tRNA synthetase specific to the amino acid valine discriminates against the very similar side chain of the amino acid isoleucine.[50]
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Proteins can bind to other proteins as well as to small-molecule substrates. When proteins bind specifically to other copies of the same molecule, they can oligomerize to form fibrils; this process occurs often in structural proteins that consist of globular monomers that self-associate to form rigid fibers. Protein–protein interactions also regulate enzymatic activity, control progression through the cell cycle, and allow the assembly of large protein complexes that carry out many closely related reactions with a common biological function. Proteins can also bind to, or even be integrated into, cell membranes. The ability of binding partners to induce conformational changes in proteins allows the construction of enormously complex signaling networks.[51]
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As interactions between proteins are reversible, and depend heavily on the availability of different groups of partner proteins to form aggregates that are capable to carry out discrete sets of function, study of the interactions between specific proteins is a key to understand important aspects of cellular function, and ultimately the properties that distinguish particular cell types.[52][53]
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The best-known role of proteins in the cell is as enzymes, which catalyse chemical reactions. Enzymes are usually highly specific and accelerate only one or a few chemical reactions. Enzymes carry out most of the reactions involved in metabolism, as well as manipulating DNA in processes such as DNA replication, DNA repair, and transcription. Some enzymes act on other proteins to add or remove chemical groups in a process known as posttranslational modification. About 4,000 reactions are known to be catalysed by enzymes.[54] The rate acceleration conferred by enzymatic catalysis is often enormous—as much as 1017-fold increase in rate over the uncatalysed reaction in the case of orotate decarboxylase (78 million years without the enzyme, 18 milliseconds with the enzyme).[55]
|
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The molecules bound and acted upon by enzymes are called substrates. Although enzymes can consist of hundreds of amino acids, it is usually only a small fraction of the residues that come in contact with the substrate, and an even smaller fraction—three to four residues on average—that are directly involved in catalysis.[56] The region of the enzyme that binds the substrate and contains the catalytic residues is known as the active site.
|
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Dirigent proteins are members of a class of proteins that dictate the stereochemistry of a compound synthesized by other enzymes.[57]
|
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Many proteins are involved in the process of cell signaling and signal transduction. Some proteins, such as insulin, are extracellular proteins that transmit a signal from the cell in which they were synthesized to other cells in distant tissues. Others are membrane proteins that act as receptors whose main function is to bind a signaling molecule and induce a biochemical response in the cell. Many receptors have a binding site exposed on the cell surface and an effector domain within the cell, which may have enzymatic activity or may undergo a conformational change detected by other proteins within the cell.[58]
|
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Antibodies are protein components of an adaptive immune system whose main function is to bind antigens, or foreign substances in the body, and target them for destruction. Antibodies can be secreted into the extracellular environment or anchored in the membranes of specialized B cells known as plasma cells. Whereas enzymes are limited in their binding affinity for their substrates by the necessity of conducting their reaction, antibodies have no such constraints. An antibody's binding affinity to its target is extraordinarily high.[59]
|
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Many ligand transport proteins bind particular small biomolecules and transport them to other locations in the body of a multicellular organism. These proteins must have a high binding affinity when their ligand is present in high concentrations, but must also release the ligand when it is present at low concentrations in the target tissues. The canonical example of a ligand-binding protein is haemoglobin, which transports oxygen from the lungs to other organs and tissues in all vertebrates and has close homologs in every biological kingdom.[60] Lectins are sugar-binding proteins which are highly specific for their sugar moieties. Lectins typically play a role in biological recognition phenomena involving cells and proteins.[61] Receptors and hormones are highly specific binding proteins.
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Transmembrane proteins can also serve as ligand transport proteins that alter the permeability of the cell membrane to small molecules and ions. The membrane alone has a hydrophobic core through which polar or charged molecules cannot diffuse. Membrane proteins contain internal channels that allow such molecules to enter and exit the cell. Many ion channel proteins are specialized to select for only a particular ion; for example, potassium and sodium channels often discriminate for only one of the two ions.[62]
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Structural proteins confer stiffness and rigidity to otherwise-fluid biological components. Most structural proteins are fibrous proteins; for example, collagen and elastin are critical components of connective tissue such as cartilage, and keratin is found in hard or filamentous structures such as hair, nails, feathers, hooves, and some animal shells.[63] Some globular proteins can also play structural functions, for example, actin and tubulin are globular and soluble as monomers, but polymerize to form long, stiff fibers that make up the cytoskeleton, which allows the cell to maintain its shape and size.
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Other proteins that serve structural functions are motor proteins such as myosin, kinesin, and dynein, which are capable of generating mechanical forces. These proteins are crucial for cellular motility of single celled organisms and the sperm of many multicellular organisms which reproduce sexually. They also generate the forces exerted by contracting muscles[64] and play essential roles in intracellular transport.
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The activities and structures of proteins may be examined in vitro, in vivo, and in silico. In vitro studies of purified proteins in controlled environments are useful for learning how a protein carries out its function: for example, enzyme kinetics studies explore the chemical mechanism of an enzyme's catalytic activity and its relative affinity for various possible substrate molecules. By contrast, in vivo experiments can provide information about the physiological role of a protein in the context of a cell or even a whole organism. In silico studies use computational methods to study proteins.
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To perform in vitro analysis, a protein must be purified away from other cellular components. This process usually begins with cell lysis, in which a cell's membrane is disrupted and its internal contents released into a solution known as a crude lysate. The resulting mixture can be purified using ultracentrifugation, which fractionates the various cellular components into fractions containing soluble proteins; membrane lipids and proteins; cellular organelles, and nucleic acids. Precipitation by a method known as salting out can concentrate the proteins from this lysate. Various types of chromatography are then used to isolate the protein or proteins of interest based on properties such as molecular weight, net charge and binding affinity.[65] The level of purification can be monitored using various types of gel electrophoresis if the desired protein's molecular weight and isoelectric point are known, by spectroscopy if the protein has distinguishable spectroscopic features, or by enzyme assays if the protein has enzymatic activity. Additionally, proteins can be isolated according to their charge using electrofocusing.[66]
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For natural proteins, a series of purification steps may be necessary to obtain protein sufficiently pure for laboratory applications. To simplify this process, genetic engineering is often used to add chemical features to proteins that make them easier to purify without affecting their structure or activity. Here, a "tag" consisting of a specific amino acid sequence, often a series of histidine residues (a "His-tag"), is attached to one terminus of the protein. As a result, when the lysate is passed over a chromatography column containing nickel, the histidine residues ligate the nickel and attach to the column while the untagged components of the lysate pass unimpeded. A number of different tags have been developed to help researchers purify specific proteins from complex mixtures.[67]
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The study of proteins in vivo is often concerned with the synthesis and localization of the protein within the cell. Although many intracellular proteins are synthesized in the cytoplasm and membrane-bound or secreted proteins in the endoplasmic reticulum, the specifics of how proteins are targeted to specific organelles or cellular structures is often unclear. A useful technique for assessing cellular localization uses genetic engineering to express in a cell a fusion protein or chimera consisting of the natural protein of interest linked to a "reporter" such as green fluorescent protein (GFP).[68] The fused protein's position within the cell can be cleanly and efficiently visualized using microscopy,[69] as shown in the figure opposite.
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Other methods for elucidating the cellular location of proteins requires the use of known compartmental markers for regions such as the ER, the Golgi, lysosomes or vacuoles, mitochondria, chloroplasts, plasma membrane, etc. With the use of fluorescently tagged versions of these markers or of antibodies to known markers, it becomes much simpler to identify the localization of a protein of interest. For example, indirect immunofluorescence will allow for fluorescence colocalization and demonstration of location. Fluorescent dyes are used to label cellular compartments for a similar purpose.[70]
|
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Other possibilities exist, as well. For example, immunohistochemistry usually utilizes an antibody to one or more proteins of interest that are conjugated to enzymes yielding either luminescent or chromogenic signals that can be compared between samples, allowing for localization information. Another applicable technique is cofractionation in sucrose (or other material) gradients using isopycnic centrifugation.[71] While this technique does not prove colocalization of a compartment of known density and the protein of interest, it does increase the likelihood, and is more amenable to large-scale studies.
|
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Finally, the gold-standard method of cellular localization is immunoelectron microscopy. This technique also uses an antibody to the protein of interest, along with classical electron microscopy techniques. The sample is prepared for normal electron microscopic examination, and then treated with an antibody to the protein of interest that is conjugated to an extremely electro-dense material, usually gold. This allows for the localization of both ultrastructural details as well as the protein of interest.[72]
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Through another genetic engineering application known as site-directed mutagenesis, researchers can alter the protein sequence and hence its structure, cellular localization, and susceptibility to regulation. This technique even allows the incorporation of unnatural amino acids into proteins, using modified tRNAs,[73] and may allow the rational design of new proteins with novel properties.[74]
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The total complement of proteins present at a time in a cell or cell type is known as its proteome, and the study of such large-scale data sets defines the field of proteomics, named by analogy to the related field of genomics. Key experimental techniques in proteomics include 2D electrophoresis,[75] which allows the separation of many proteins, mass spectrometry,[76] which allows rapid high-throughput identification of proteins and sequencing of peptides (most often after in-gel digestion), protein microarrays, which allow the detection of the relative levels of the various proteins present in a cell, and two-hybrid screening, which allows the systematic exploration of protein–protein interactions.[77] The total complement of biologically possible such interactions is known as the interactome.[78] A systematic attempt to determine the structures of proteins representing every possible fold is known as structural genomics.[79]
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A vast array of computational methods have been developed to analyze the structure, function and evolution of proteins.
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The development of such tools has been driven by the large amount of genomic and proteomic data available for a variety of organisms, including the human genome. It is simply impossible to study all proteins experimentally, hence only a few are subjected to laboratory experiments while computational tools are used to extrapolate to similar proteins. Such homologous proteins can be efficiently identified in distantly related organisms by sequence alignment. Genome and gene sequences can be searched by a variety of tools for certain properties. Sequence profiling tools can find restriction enzyme sites, open reading frames in nucleotide sequences, and predict secondary structures. Phylogenetic trees can be constructed and evolutionary hypotheses developed using special software like ClustalW regarding the ancestry of modern organisms and the genes they express. The field of bioinformatics is now indispensable for the analysis of genes and proteins.
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Discovering the tertiary structure of a protein, or the quaternary structure of its complexes, can provide important clues about how the protein performs its function and how it can be affected, i.e. in drug design. As proteins are too small to be seen under a light microscope, other methods have to be employed to determine their structure. Common experimental methods include X-ray crystallography and NMR spectroscopy, both of which can produce structural information at atomic resolution. However, NMR experiments are able to provide information from which a subset of distances between pairs of atoms can be estimated, and the final possible conformations for a protein are determined by solving a distance geometry problem. Dual polarisation interferometry is a quantitative analytical method for measuring the overall protein conformation and conformational changes due to interactions or other stimulus. Circular dichroism is another laboratory technique for determining internal β-sheet / α-helical composition of proteins. Cryoelectron microscopy is used to produce lower-resolution structural information about very large protein complexes, including assembled viruses;[80] a variant known as electron crystallography can also produce high-resolution information in some cases, especially for two-dimensional crystals of membrane proteins.[81] Solved structures are usually deposited in the Protein Data Bank (PDB), a freely available resource from which structural data about thousands of proteins can be obtained in the form of Cartesian coordinates for each atom in the protein.[82]
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Many more gene sequences are known than protein structures. Further, the set of solved structures is biased toward proteins that can be easily subjected to the conditions required in X-ray crystallography, one of the major structure determination methods. In particular, globular proteins are comparatively easy to crystallize in preparation for X-ray crystallography. Membrane proteins and large protein complexes, by contrast, are difficult to crystallize and are underrepresented in the PDB.[83] Structural genomics initiatives have attempted to remedy these deficiencies by systematically solving representative structures of major fold classes. Protein structure prediction methods attempt to provide a means of generating a plausible structure for proteins whose structures have not been experimentally determined.[84]
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Complementary to the field of structural genomics, protein structure prediction develops efficient mathematical models of proteins to computationally predict the molecular formations in theory, instead of detecting structures with laboratory observation.[85] The most successful type of structure prediction, known as homology modeling, relies on the existence of a "template" structure with sequence similarity to the protein being modeled; structural genomics' goal is to provide sufficient representation in solved structures to model most of those that remain.[86] Although producing accurate models remains a challenge when only distantly related template structures are available, it has been suggested that sequence alignment is the bottleneck in this process, as quite accurate models can be produced if a "perfect" sequence alignment is known.[87] Many structure prediction methods have served to inform the emerging field of protein engineering, in which novel protein folds have already been designed.[88] A more complex computational problem is the prediction of intermolecular interactions, such as in molecular docking and protein–protein interaction prediction.[89]
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Mathematical models to simulate dynamic processes of protein folding and binding involve molecular mechanics, in particular, molecular dynamics. Monte Carlo techniques facilitate the computations, which exploit advances in parallel and distributed computing (for example, the Folding@home project[90] which performs molecular modeling on GPUs). In silico simulations discovered the folding of small α-helical protein domains such as the villin headpiece[91] and the HIV accessory protein.[92] Hybrid methods combining standard molecular dynamics with quantum mechanical mathematics explored the electronic states of rhodopsins.[93]
|
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Many proteins (in Eucaryota ~33%) contain large unstructured but biologically functional segments and can be classified as intrinsically disordered proteins.[94] Predicting and analysing protein disorder is, therefore, an important part of protein structure characterisation.[95]
|
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Most microorganisms and plants can biosynthesize all 20 standard amino acids, while animals (including humans) must obtain some of the amino acids from the diet.[49] The amino acids that an organism cannot synthesize on its own are referred to as essential amino acids. Key enzymes that synthesize certain amino acids are not present in animals—such as aspartokinase, which catalyses the first step in the synthesis of lysine, methionine, and threonine from aspartate. If amino acids are present in the environment, microorganisms can conserve energy by taking up the amino acids from their surroundings and downregulating their biosynthetic pathways.
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In animals, amino acids are obtained through the consumption of foods containing protein. Ingested proteins are then broken down into amino acids through digestion, which typically involves denaturation of the protein through exposure to acid and hydrolysis by enzymes called proteases. Some ingested amino acids are used for protein biosynthesis, while others are converted to glucose through gluconeogenesis, or fed into the citric acid cycle. This use of protein as a fuel is particularly important under starvation conditions as it allows the body's own proteins to be used to support life, particularly those found in muscle.[96]
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In animals such as dogs and cats, protein maintains the health and quality of the skin by promoting hair follicle growth and keratinization, and thus reducing the likelihood of skin problems producing malodours.[97] Poor-quality proteins also have a role regarding gastrointestinal health, increasing the potential for flatulence and odorous compounds in dogs because when proteins reach the colon in an undigested state, they are fermented producing hydrogen sulfide gas, indole, and skatole.[98] Dogs and cats digest animal proteins better than those from plants, but products of low-quality animal origin are poorly digested, including skin, feathers, and connective tissue.[98]
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Proteins are large biomolecules, or macromolecules, consisting of one or more long chains of amino acid residues. Proteins perform a vast array of functions within organisms, including catalysing metabolic reactions, DNA replication, responding to stimuli, providing structure to cells, and organisms, and transporting molecules from one location to another. Proteins differ from one another primarily in their sequence of amino acids, which is dictated by the nucleotide sequence of their genes, and which usually results in protein folding into a specific 3D structure that determines its activity.
|
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A linear chain of amino acid residues is called a polypeptide. A protein contains at least one long polypeptide. Short polypeptides, containing less than 20–30 residues, are rarely considered to be proteins and are commonly called peptides, or sometimes oligopeptides. The individual amino acid residues are bonded together by peptide bonds and adjacent amino acid residues. The sequence of amino acid residues in a protein is defined by the sequence of a gene, which is encoded in the genetic code. In general, the genetic code specifies 20 standard amino acids; but in certain organisms the genetic code can include selenocysteine and—in certain archaea—pyrrolysine. Shortly after or even during synthesis, the residues in a protein are often chemically modified by post-translational modification, which alters the physical and chemical properties, folding, stability, activity, and ultimately, the function of the proteins. Some proteins have non-peptide groups attached, which can be called prosthetic groups or cofactors. Proteins can also work together to achieve a particular function, and they often associate to form stable protein complexes.
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Once formed, proteins only exist for a certain period and are then degraded and recycled by the cell's machinery through the process of protein turnover. A protein's lifespan is measured in terms of its half-life and covers a wide range. They can exist for minutes or years with an average lifespan of 1–2 days in mammalian cells. Abnormal or misfolded proteins are degraded more rapidly either due to being targeted for destruction or due to being unstable.
|
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Like other biological macromolecules such as polysaccharides and nucleic acids, proteins are essential parts of organisms and participate in virtually every process within cells. Many proteins are enzymes that catalyse biochemical reactions and are vital to metabolism. Proteins also have structural or mechanical functions, such as actin and myosin in muscle and the proteins in the cytoskeleton, which form a system of scaffolding that maintains cell shape. Other proteins are important in cell signaling, immune responses, cell adhesion, and the cell cycle. In animals, proteins are needed in the diet to provide the essential amino acids that cannot be synthesized. Digestion breaks the proteins down for use in the metabolism.
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Proteins may be purified from other cellular components using a variety of techniques such as ultracentrifugation, precipitation, electrophoresis, and chromatography; the advent of genetic engineering has made possible a number of methods to facilitate purification. Methods commonly used to study protein structure and function include immunohistochemistry, site-directed mutagenesis, X-ray crystallography, nuclear magnetic resonance and mass spectrometry.
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Proteins were recognized as a distinct class of biological molecules in the eighteenth century by Antoine Fourcroy and others, distinguished by the molecules' ability to coagulate or flocculate under treatments with heat or acid.[1] Noted examples at the time included albumin from egg whites, blood serum albumin, fibrin, and wheat gluten.
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Proteins were first described by the Dutch chemist Gerardus Johannes Mulder and named by the Swedish chemist Jöns Jacob Berzelius in 1838.[2][3] Mulder carried out elemental analysis of common proteins and found that nearly all proteins had the same empirical formula, C400H620N100O120P1S1.[4] He came to the erroneous conclusion that they might be composed of a single type of (very large) molecule. The term "protein" to describe these molecules was proposed by Mulder's associate Berzelius; protein is derived from the Greek word πρώτειος (proteios), meaning "primary",[5] "in the lead", or "standing in front",[6] + -in. Mulder went on to identify the products of protein degradation such as the amino acid leucine for which he found a (nearly correct) molecular weight of 131 Da.[4] Prior to "protein", other names were used, like "albumins" or "albuminous materials" (Eiweisskörper, in German).[7]
|
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Early nutritional scientists such as the German Carl von Voit believed that protein was the most important nutrient for maintaining the structure of the body, because it was generally believed that "flesh makes flesh."[8] Karl Heinrich Ritthausen extended known protein forms with the identification of glutamic acid. At the Connecticut Agricultural Experiment Station a detailed review of the vegetable proteins was compiled by Thomas Burr Osborne. Working with Lafayette Mendel and applying Liebig's law of the minimum in feeding laboratory rats, the nutritionally essential amino acids were established. The work was continued and communicated by William Cumming Rose. The understanding of proteins as polypeptides came through the work of Franz Hofmeister and Hermann Emil Fischer in 1902.[9][10] The central role of proteins as enzymes in living organisms was not fully appreciated until 1926, when James B. Sumner showed that the enzyme urease was in fact a protein.[11]
|
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The difficulty in purifying proteins in large quantities made them very difficult for early protein biochemists to study. Hence, early studies focused on proteins that could be purified in large quantities, e.g., those of blood, egg white, various toxins, and digestive/metabolic enzymes obtained from slaughterhouses. In the 1950s, the Armour Hot Dog Co. purified 1 kg of pure bovine pancreatic ribonuclease A and made it freely available to scientists; this gesture helped ribonuclease A become a major target for biochemical study for the following decades.[4]
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Linus Pauling is credited with the successful prediction of regular protein secondary structures based on hydrogen bonding, an idea first put forth by William Astbury in 1933.[12] Later work by Walter Kauzmann on denaturation,[13][14] based partly on previous studies by Kaj Linderstrøm-Lang,[15] contributed an understanding of protein folding and structure mediated by hydrophobic interactions.
|
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The first protein to be sequenced was insulin, by Frederick Sanger, in 1949. Sanger correctly determined the amino acid sequence of insulin, thus conclusively demonstrating that proteins consisted of linear polymers of amino acids rather than branched chains, colloids, or cyclols.[16] He won the Nobel Prize for this achievement in 1958.[17]
|
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The first protein structures to be solved were hemoglobin and myoglobin, by Max Perutz and Sir John Cowdery Kendrew, respectively, in 1958.[18][19] As of 2017[update], the Protein Data Bank has over 126,060 atomic-resolution structures of proteins.[20] In more recent times, cryo-electron microscopy of large macromolecular assemblies[21] and computational protein structure prediction of small protein domains[22] are two methods approaching atomic resolution.
|
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Most proteins consist of linear polymers built from series of up to 20 different L-α- amino acids. All proteinogenic amino acids possess common structural features, including an α-carbon to which an amino group, a carboxyl group, and a variable side chain are bonded. Only proline differs from this basic structure as it contains an unusual ring to the N-end amine group, which forces the CO–NH amide moiety into a fixed conformation.[23] The side chains of the standard amino acids, detailed in the list of standard amino acids, have a great variety of chemical structures and properties; it is the combined effect of all of the amino acid side chains in a protein that ultimately determines its three-dimensional structure and its chemical reactivity.[24]
|
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The amino acids in a polypeptide chain are linked by peptide bonds. Once linked in the protein chain, an individual amino acid is called a residue, and the linked series of carbon, nitrogen, and oxygen atoms are known as the main chain or protein backbone.[25]
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The peptide bond has two resonance forms that contribute some double-bond character and inhibit rotation around its axis, so that the alpha carbons are roughly coplanar. The other two dihedral angles in the peptide bond determine the local shape assumed by the protein backbone.[26] The end with a free amino group is known as the N-terminus or amino terminus, whereas the end of the protein with a free carboxyl group is known as the C-terminus or carboxy terminus (the sequence of the protein is written from N-terminus to C-terminus, from left to right).
|
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The words protein, polypeptide, and peptide are a little ambiguous and can overlap in meaning. Protein is generally used to refer to the complete biological molecule in a stable conformation, whereas peptide is generally reserved for a short amino acid oligomers often lacking a stable 3D structure. But the boundary between the two is not well defined and usually lies near 20–30 residues.[27] Polypeptide can refer to any single linear chain of amino acids, usually regardless of length, but often implies an absence of a defined conformation.
|
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+
Proteins can interact with many types of molecules, including with other proteins, with lipids, with carboyhydrates, and with DNA.[28][29][30][31]
|
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+
It has been estimated that average-sized bacteria contain about 2 million proteins per cell (e.g. E. coli and Staphylococcus aureus). Smaller bacteria, such as Mycoplasma or spirochetes contain fewer molecules, on the order of 50,000 to 1 million. By contrast, eukaryotic cells are larger and thus contain much more protein. For instance, yeast cells have been estimated to contain about 50 million proteins and human cells on the order of 1 to 3 billion.[32] The concentration of individual protein copies ranges from a few molecules per cell up to 20 million.[33] Not all genes coding proteins are expressed in most cells and their number depends on, for example, cell type and external stimuli. For instance, of the 20,000 or so proteins encoded by the human genome, only 6,000 are detected in lymphoblastoid cells.[34] Moreover, the number of proteins the genome encodes correlates well with the organism complexity. Eukaryotes have 15,000, bacteria have 3,200, archaea have 2,400, and viruses have 42 proteins on average coded in their respective genomes.[35]
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Proteins are assembled from amino acids using information encoded in genes. Each protein has its own unique amino acid sequence that is specified by the nucleotide sequence of the gene encoding this protein. The genetic code is a set of three-nucleotide sets called codons and each three-nucleotide combination designates an amino acid, for example AUG (adenine–uracil–guanine) is the code for methionine. Because DNA contains four nucleotides, the total number of possible codons is 64; hence, there is some redundancy in the genetic code, with some amino acids specified by more than one codon.[36] Genes encoded in DNA are first transcribed into pre-messenger RNA (mRNA) by proteins such as RNA polymerase. Most organisms then process the pre-mRNA (also known as a primary transcript) using various forms of Post-transcriptional modification to form the mature mRNA, which is then used as a template for protein synthesis by the ribosome. In prokaryotes the mRNA may either be used as soon as it is produced, or be bound by a ribosome after having moved away from the nucleoid. In contrast, eukaryotes make mRNA in the cell nucleus and then translocate it across the nuclear membrane into the cytoplasm, where protein synthesis then takes place. The rate of protein synthesis is higher in prokaryotes than eukaryotes and can reach up to 20 amino acids per second.[37]
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The process of synthesizing a protein from an mRNA template is known as translation. The mRNA is loaded onto the ribosome and is read three nucleotides at a time by matching each codon to its base pairing anticodon located on a transfer RNA molecule, which carries the amino acid corresponding to the codon it recognizes. The enzyme aminoacyl tRNA synthetase "charges" the tRNA molecules with the correct amino acids. The growing polypeptide is often termed the nascent chain. Proteins are always biosynthesized from N-terminus to C-terminus.[36]
|
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|
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+
The size of a synthesized protein can be measured by the number of amino acids it contains and by its total molecular mass, which is normally reported in units of daltons (synonymous with atomic mass units), or the derivative unit kilodalton (kDa). The average size of a protein increases from Archaea to Bacteria to Eukaryote (283, 311, 438 residues and 31, 34, 49 kDa respectively) due to a bigger number of protein domains constituting proteins in higher organisms.[35] For instance, yeast proteins are on average 466 amino acids long and 53 kDa in mass.[27] The largest known proteins are the titins, a component of the muscle sarcomere, with a molecular mass of almost 3,000 kDa and a total length of almost 27,000 amino acids.[38]
|
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+
|
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+
Short proteins can also be synthesized chemically by a family of methods known as peptide synthesis, which rely on organic synthesis techniques such as chemical ligation to produce peptides in high yield.[39] Chemical synthesis allows for the introduction of non-natural amino acids into polypeptide chains, such as attachment of fluorescent probes to amino acid side chains.[40] These methods are useful in laboratory biochemistry and cell biology, though generally not for commercial applications. Chemical synthesis is inefficient for polypeptides longer than about 300 amino acids, and the synthesized proteins may not readily assume their native tertiary structure. Most chemical synthesis methods proceed from C-terminus to N-terminus, opposite the biological reaction.[41]
|
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+
|
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+
Most proteins fold into unique 3D structures. The shape into which a protein naturally folds is known as its native conformation.[42] Although many proteins can fold unassisted, simply through the chemical properties of their amino acids, others require the aid of molecular chaperones to fold into their native states.[43] Biochemists often refer to four distinct aspects of a protein's structure:[44]
|
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+
|
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+
Proteins are not entirely rigid molecules. In addition to these levels of structure, proteins may shift between several related structures while they perform their functions. In the context of these functional rearrangements, these tertiary or quaternary structures are usually referred to as "conformations", and transitions between them are called conformational changes. Such changes are often induced by the binding of a substrate molecule to an enzyme's active site, or the physical region of the protein that participates in chemical catalysis. In solution proteins also undergo variation in structure through thermal vibration and the collision with other molecules.[45]
|
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|
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+
Proteins can be informally divided into three main classes, which correlate with typical tertiary structures: globular proteins, fibrous proteins, and membrane proteins. Almost all globular proteins are soluble and many are enzymes. Fibrous proteins are often structural, such as collagen, the major component of connective tissue, or keratin, the protein component of hair and nails. Membrane proteins often serve as receptors or provide channels for polar or charged molecules to pass through the cell membrane.[46]
|
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+
|
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+
A special case of intramolecular hydrogen bonds within proteins, poorly shielded from water attack and hence promoting their own dehydration, are called dehydrons.[47]
|
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+
|
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+
Many proteins are composed of several protein domains, i.e. segments of a protein that fold into distinct structural units. Domains usually also have specific functions, such as enzymatic activities (e.g. kinase) or they serve as binding modules (e.g. the SH3 domain binds to proline-rich sequences in other proteins).
|
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+
|
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+
Short amino acid sequences within proteins often act as recognition sites for other proteins.[48] For instance, SH3 domains typically bind to short PxxP motifs (i.e. 2 prolines [P], separated by two unspecified amino acids [x], although the surrounding amino acids may determine the exact binding specificity). Many such motifs has been collected in the Eukaryotic Linear Motif (ELM) database.
|
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Proteins are the chief actors within the cell, said to be carrying out the duties specified by the information encoded in genes.[27] With the exception of certain types of RNA, most other biological molecules are relatively inert elements upon which proteins act. Proteins make up half the dry weight of an Escherichia coli cell, whereas other macromolecules such as DNA and RNA make up only 3% and 20%, respectively.[49] The set of proteins expressed in a particular cell or cell type is known as its proteome.
|
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+
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+
The chief characteristic of proteins that also allows their diverse set of functions is their ability to bind other molecules specifically and tightly. The region of the protein responsible for binding another molecule is known as the binding site and is often a depression or "pocket" on the molecular surface. This binding ability is mediated by the tertiary structure of the protein, which defines the binding site pocket, and by the chemical properties of the surrounding amino acids' side chains. Protein binding can be extraordinarily tight and specific; for example, the ribonuclease inhibitor protein binds to human angiogenin with a sub-femtomolar dissociation constant (<10−15 M) but does not bind at all to its amphibian homolog onconase (>1 M). Extremely minor chemical changes such as the addition of a single methyl group to a binding partner can sometimes suffice to nearly eliminate binding; for example, the aminoacyl tRNA synthetase specific to the amino acid valine discriminates against the very similar side chain of the amino acid isoleucine.[50]
|
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Proteins can bind to other proteins as well as to small-molecule substrates. When proteins bind specifically to other copies of the same molecule, they can oligomerize to form fibrils; this process occurs often in structural proteins that consist of globular monomers that self-associate to form rigid fibers. Protein–protein interactions also regulate enzymatic activity, control progression through the cell cycle, and allow the assembly of large protein complexes that carry out many closely related reactions with a common biological function. Proteins can also bind to, or even be integrated into, cell membranes. The ability of binding partners to induce conformational changes in proteins allows the construction of enormously complex signaling networks.[51]
|
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As interactions between proteins are reversible, and depend heavily on the availability of different groups of partner proteins to form aggregates that are capable to carry out discrete sets of function, study of the interactions between specific proteins is a key to understand important aspects of cellular function, and ultimately the properties that distinguish particular cell types.[52][53]
|
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The best-known role of proteins in the cell is as enzymes, which catalyse chemical reactions. Enzymes are usually highly specific and accelerate only one or a few chemical reactions. Enzymes carry out most of the reactions involved in metabolism, as well as manipulating DNA in processes such as DNA replication, DNA repair, and transcription. Some enzymes act on other proteins to add or remove chemical groups in a process known as posttranslational modification. About 4,000 reactions are known to be catalysed by enzymes.[54] The rate acceleration conferred by enzymatic catalysis is often enormous—as much as 1017-fold increase in rate over the uncatalysed reaction in the case of orotate decarboxylase (78 million years without the enzyme, 18 milliseconds with the enzyme).[55]
|
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|
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The molecules bound and acted upon by enzymes are called substrates. Although enzymes can consist of hundreds of amino acids, it is usually only a small fraction of the residues that come in contact with the substrate, and an even smaller fraction—three to four residues on average—that are directly involved in catalysis.[56] The region of the enzyme that binds the substrate and contains the catalytic residues is known as the active site.
|
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|
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Dirigent proteins are members of a class of proteins that dictate the stereochemistry of a compound synthesized by other enzymes.[57]
|
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Many proteins are involved in the process of cell signaling and signal transduction. Some proteins, such as insulin, are extracellular proteins that transmit a signal from the cell in which they were synthesized to other cells in distant tissues. Others are membrane proteins that act as receptors whose main function is to bind a signaling molecule and induce a biochemical response in the cell. Many receptors have a binding site exposed on the cell surface and an effector domain within the cell, which may have enzymatic activity or may undergo a conformational change detected by other proteins within the cell.[58]
|
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|
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Antibodies are protein components of an adaptive immune system whose main function is to bind antigens, or foreign substances in the body, and target them for destruction. Antibodies can be secreted into the extracellular environment or anchored in the membranes of specialized B cells known as plasma cells. Whereas enzymes are limited in their binding affinity for their substrates by the necessity of conducting their reaction, antibodies have no such constraints. An antibody's binding affinity to its target is extraordinarily high.[59]
|
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Many ligand transport proteins bind particular small biomolecules and transport them to other locations in the body of a multicellular organism. These proteins must have a high binding affinity when their ligand is present in high concentrations, but must also release the ligand when it is present at low concentrations in the target tissues. The canonical example of a ligand-binding protein is haemoglobin, which transports oxygen from the lungs to other organs and tissues in all vertebrates and has close homologs in every biological kingdom.[60] Lectins are sugar-binding proteins which are highly specific for their sugar moieties. Lectins typically play a role in biological recognition phenomena involving cells and proteins.[61] Receptors and hormones are highly specific binding proteins.
|
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Transmembrane proteins can also serve as ligand transport proteins that alter the permeability of the cell membrane to small molecules and ions. The membrane alone has a hydrophobic core through which polar or charged molecules cannot diffuse. Membrane proteins contain internal channels that allow such molecules to enter and exit the cell. Many ion channel proteins are specialized to select for only a particular ion; for example, potassium and sodium channels often discriminate for only one of the two ions.[62]
|
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Structural proteins confer stiffness and rigidity to otherwise-fluid biological components. Most structural proteins are fibrous proteins; for example, collagen and elastin are critical components of connective tissue such as cartilage, and keratin is found in hard or filamentous structures such as hair, nails, feathers, hooves, and some animal shells.[63] Some globular proteins can also play structural functions, for example, actin and tubulin are globular and soluble as monomers, but polymerize to form long, stiff fibers that make up the cytoskeleton, which allows the cell to maintain its shape and size.
|
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Other proteins that serve structural functions are motor proteins such as myosin, kinesin, and dynein, which are capable of generating mechanical forces. These proteins are crucial for cellular motility of single celled organisms and the sperm of many multicellular organisms which reproduce sexually. They also generate the forces exerted by contracting muscles[64] and play essential roles in intracellular transport.
|
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The activities and structures of proteins may be examined in vitro, in vivo, and in silico. In vitro studies of purified proteins in controlled environments are useful for learning how a protein carries out its function: for example, enzyme kinetics studies explore the chemical mechanism of an enzyme's catalytic activity and its relative affinity for various possible substrate molecules. By contrast, in vivo experiments can provide information about the physiological role of a protein in the context of a cell or even a whole organism. In silico studies use computational methods to study proteins.
|
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|
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To perform in vitro analysis, a protein must be purified away from other cellular components. This process usually begins with cell lysis, in which a cell's membrane is disrupted and its internal contents released into a solution known as a crude lysate. The resulting mixture can be purified using ultracentrifugation, which fractionates the various cellular components into fractions containing soluble proteins; membrane lipids and proteins; cellular organelles, and nucleic acids. Precipitation by a method known as salting out can concentrate the proteins from this lysate. Various types of chromatography are then used to isolate the protein or proteins of interest based on properties such as molecular weight, net charge and binding affinity.[65] The level of purification can be monitored using various types of gel electrophoresis if the desired protein's molecular weight and isoelectric point are known, by spectroscopy if the protein has distinguishable spectroscopic features, or by enzyme assays if the protein has enzymatic activity. Additionally, proteins can be isolated according to their charge using electrofocusing.[66]
|
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|
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For natural proteins, a series of purification steps may be necessary to obtain protein sufficiently pure for laboratory applications. To simplify this process, genetic engineering is often used to add chemical features to proteins that make them easier to purify without affecting their structure or activity. Here, a "tag" consisting of a specific amino acid sequence, often a series of histidine residues (a "His-tag"), is attached to one terminus of the protein. As a result, when the lysate is passed over a chromatography column containing nickel, the histidine residues ligate the nickel and attach to the column while the untagged components of the lysate pass unimpeded. A number of different tags have been developed to help researchers purify specific proteins from complex mixtures.[67]
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The study of proteins in vivo is often concerned with the synthesis and localization of the protein within the cell. Although many intracellular proteins are synthesized in the cytoplasm and membrane-bound or secreted proteins in the endoplasmic reticulum, the specifics of how proteins are targeted to specific organelles or cellular structures is often unclear. A useful technique for assessing cellular localization uses genetic engineering to express in a cell a fusion protein or chimera consisting of the natural protein of interest linked to a "reporter" such as green fluorescent protein (GFP).[68] The fused protein's position within the cell can be cleanly and efficiently visualized using microscopy,[69] as shown in the figure opposite.
|
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|
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Other methods for elucidating the cellular location of proteins requires the use of known compartmental markers for regions such as the ER, the Golgi, lysosomes or vacuoles, mitochondria, chloroplasts, plasma membrane, etc. With the use of fluorescently tagged versions of these markers or of antibodies to known markers, it becomes much simpler to identify the localization of a protein of interest. For example, indirect immunofluorescence will allow for fluorescence colocalization and demonstration of location. Fluorescent dyes are used to label cellular compartments for a similar purpose.[70]
|
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+
|
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+
Other possibilities exist, as well. For example, immunohistochemistry usually utilizes an antibody to one or more proteins of interest that are conjugated to enzymes yielding either luminescent or chromogenic signals that can be compared between samples, allowing for localization information. Another applicable technique is cofractionation in sucrose (or other material) gradients using isopycnic centrifugation.[71] While this technique does not prove colocalization of a compartment of known density and the protein of interest, it does increase the likelihood, and is more amenable to large-scale studies.
|
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|
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+
Finally, the gold-standard method of cellular localization is immunoelectron microscopy. This technique also uses an antibody to the protein of interest, along with classical electron microscopy techniques. The sample is prepared for normal electron microscopic examination, and then treated with an antibody to the protein of interest that is conjugated to an extremely electro-dense material, usually gold. This allows for the localization of both ultrastructural details as well as the protein of interest.[72]
|
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Through another genetic engineering application known as site-directed mutagenesis, researchers can alter the protein sequence and hence its structure, cellular localization, and susceptibility to regulation. This technique even allows the incorporation of unnatural amino acids into proteins, using modified tRNAs,[73] and may allow the rational design of new proteins with novel properties.[74]
|
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The total complement of proteins present at a time in a cell or cell type is known as its proteome, and the study of such large-scale data sets defines the field of proteomics, named by analogy to the related field of genomics. Key experimental techniques in proteomics include 2D electrophoresis,[75] which allows the separation of many proteins, mass spectrometry,[76] which allows rapid high-throughput identification of proteins and sequencing of peptides (most often after in-gel digestion), protein microarrays, which allow the detection of the relative levels of the various proteins present in a cell, and two-hybrid screening, which allows the systematic exploration of protein–protein interactions.[77] The total complement of biologically possible such interactions is known as the interactome.[78] A systematic attempt to determine the structures of proteins representing every possible fold is known as structural genomics.[79]
|
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|
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+
A vast array of computational methods have been developed to analyze the structure, function and evolution of proteins.
|
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The development of such tools has been driven by the large amount of genomic and proteomic data available for a variety of organisms, including the human genome. It is simply impossible to study all proteins experimentally, hence only a few are subjected to laboratory experiments while computational tools are used to extrapolate to similar proteins. Such homologous proteins can be efficiently identified in distantly related organisms by sequence alignment. Genome and gene sequences can be searched by a variety of tools for certain properties. Sequence profiling tools can find restriction enzyme sites, open reading frames in nucleotide sequences, and predict secondary structures. Phylogenetic trees can be constructed and evolutionary hypotheses developed using special software like ClustalW regarding the ancestry of modern organisms and the genes they express. The field of bioinformatics is now indispensable for the analysis of genes and proteins.
|
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|
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+
Discovering the tertiary structure of a protein, or the quaternary structure of its complexes, can provide important clues about how the protein performs its function and how it can be affected, i.e. in drug design. As proteins are too small to be seen under a light microscope, other methods have to be employed to determine their structure. Common experimental methods include X-ray crystallography and NMR spectroscopy, both of which can produce structural information at atomic resolution. However, NMR experiments are able to provide information from which a subset of distances between pairs of atoms can be estimated, and the final possible conformations for a protein are determined by solving a distance geometry problem. Dual polarisation interferometry is a quantitative analytical method for measuring the overall protein conformation and conformational changes due to interactions or other stimulus. Circular dichroism is another laboratory technique for determining internal β-sheet / α-helical composition of proteins. Cryoelectron microscopy is used to produce lower-resolution structural information about very large protein complexes, including assembled viruses;[80] a variant known as electron crystallography can also produce high-resolution information in some cases, especially for two-dimensional crystals of membrane proteins.[81] Solved structures are usually deposited in the Protein Data Bank (PDB), a freely available resource from which structural data about thousands of proteins can be obtained in the form of Cartesian coordinates for each atom in the protein.[82]
|
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|
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Many more gene sequences are known than protein structures. Further, the set of solved structures is biased toward proteins that can be easily subjected to the conditions required in X-ray crystallography, one of the major structure determination methods. In particular, globular proteins are comparatively easy to crystallize in preparation for X-ray crystallography. Membrane proteins and large protein complexes, by contrast, are difficult to crystallize and are underrepresented in the PDB.[83] Structural genomics initiatives have attempted to remedy these deficiencies by systematically solving representative structures of major fold classes. Protein structure prediction methods attempt to provide a means of generating a plausible structure for proteins whose structures have not been experimentally determined.[84]
|
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|
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Complementary to the field of structural genomics, protein structure prediction develops efficient mathematical models of proteins to computationally predict the molecular formations in theory, instead of detecting structures with laboratory observation.[85] The most successful type of structure prediction, known as homology modeling, relies on the existence of a "template" structure with sequence similarity to the protein being modeled; structural genomics' goal is to provide sufficient representation in solved structures to model most of those that remain.[86] Although producing accurate models remains a challenge when only distantly related template structures are available, it has been suggested that sequence alignment is the bottleneck in this process, as quite accurate models can be produced if a "perfect" sequence alignment is known.[87] Many structure prediction methods have served to inform the emerging field of protein engineering, in which novel protein folds have already been designed.[88] A more complex computational problem is the prediction of intermolecular interactions, such as in molecular docking and protein–protein interaction prediction.[89]
|
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|
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Mathematical models to simulate dynamic processes of protein folding and binding involve molecular mechanics, in particular, molecular dynamics. Monte Carlo techniques facilitate the computations, which exploit advances in parallel and distributed computing (for example, the Folding@home project[90] which performs molecular modeling on GPUs). In silico simulations discovered the folding of small α-helical protein domains such as the villin headpiece[91] and the HIV accessory protein.[92] Hybrid methods combining standard molecular dynamics with quantum mechanical mathematics explored the electronic states of rhodopsins.[93]
|
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Many proteins (in Eucaryota ~33%) contain large unstructured but biologically functional segments and can be classified as intrinsically disordered proteins.[94] Predicting and analysing protein disorder is, therefore, an important part of protein structure characterisation.[95]
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|
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Most microorganisms and plants can biosynthesize all 20 standard amino acids, while animals (including humans) must obtain some of the amino acids from the diet.[49] The amino acids that an organism cannot synthesize on its own are referred to as essential amino acids. Key enzymes that synthesize certain amino acids are not present in animals—such as aspartokinase, which catalyses the first step in the synthesis of lysine, methionine, and threonine from aspartate. If amino acids are present in the environment, microorganisms can conserve energy by taking up the amino acids from their surroundings and downregulating their biosynthetic pathways.
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In animals, amino acids are obtained through the consumption of foods containing protein. Ingested proteins are then broken down into amino acids through digestion, which typically involves denaturation of the protein through exposure to acid and hydrolysis by enzymes called proteases. Some ingested amino acids are used for protein biosynthesis, while others are converted to glucose through gluconeogenesis, or fed into the citric acid cycle. This use of protein as a fuel is particularly important under starvation conditions as it allows the body's own proteins to be used to support life, particularly those found in muscle.[96]
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In animals such as dogs and cats, protein maintains the health and quality of the skin by promoting hair follicle growth and keratinization, and thus reducing the likelihood of skin problems producing malodours.[97] Poor-quality proteins also have a role regarding gastrointestinal health, increasing the potential for flatulence and odorous compounds in dogs because when proteins reach the colon in an undigested state, they are fermented producing hydrogen sulfide gas, indole, and skatole.[98] Dogs and cats digest animal proteins better than those from plants, but products of low-quality animal origin are poorly digested, including skin, feathers, and connective tissue.[98]
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Auvergne-Rhône-Alpes (ARA; French: [ovɛʁɲ ʁon‿alp] (listen); Arpitan: Ôvèrgne-Rôno-Ârpes; Occitan: Auvèrnhe Ròse Aups; Italian: Alvernia-Rodano-Alpi) is a region in southeast-central France created by the territorial reform of French Regions in 2015; it resulted from the merger of Auvergne and Rhône-Alpes. The new region came into effect on 1 January 2016, after the regional elections in December 2015.[2]
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The region covers an area of more than 69,711 km2 (26,916 sq mi), making it the third largest in metropolitan France, and has a population of 7,877,698, second only to Île-de-France.[3] It consists of 12 departments and one territorial collectivity. Lyon is the chef-lieu of the region.
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This new region combines diverse geographical, sociological, economic, and cultural regions, which was already true of Rhône-Alpes, as well as Auvergne, to a lesser extent. While the old Rhône-Alpes and Auvergne regions each enjoyed a unity defined by axes of communication and the pull of their respective metropoles,[Note 1] the new combination is heterogeneous, and sustained lively opposition from some local officials after its creation.[4][5][6][7]
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The text of the territorial reform law gives interim names for most of the merged regions, combining the names of their constituent regions alphabetically, separated by hyphens. Permanent names would be proposed by the new regional councils and confirmed by the Conseil d'État by 1 October 2016.[8][9]
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The interim name of the new administrative region was a hyphenated placename, composed of the historic region of Auvergne, the river Rhône, and the French Alps (Alpes). The same name has been chosen as the definitive name, which was officialized by the Conseil d'État on 28 September 2016.[10]
|
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According to several online polls from Lyon Capitale, the name "Rhône-Alpes-Auvergne" led voting, ahead of "Alpes-Auvergne" and "AURA" (an acronym for Auvergne-Rhône-Alpes)[11], which was proposed by Jean-Jack Queyranne, the former president for the former Rhône-Alpes region. Schoolchildren were consulted about the name of the new region in February 2016, and local residents were consulted in March.[12]
|
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After adjusting the votes in proportion to the number of inhabitants of the regions (Rhône-Alpes having five times the population of Auvergne) the name "Rhône-Alpes-Auvergne" was still leading, ahead of "Auvergne-Rhône-Alpes" and the acronym "AURA".[13]
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Despite this result, Laurent Wauquiez and his team decided not to follow the preference of the citizens of the new region, and the name Auvergne-Rhône-Alpes was put to the vote by the regional assembly and adopted unanimously on 23 June 2016;[14][15] it was made official on 28 September 2016 through a decree appearing in the Journal officiel de la République française.[1]
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The region was given a coat of arms in October 2017, which combined those of Auvergne, Savoie, Lyonnais and Dauphiné[16] It also has a flag, which initially consisted of the coat of arms on a white background, but was replaced by a heraldic flag in January 2018. On 9 February 2018, the region formalized the flag and the coat of arms on its website, as implemented by Mattieu Casali, a historical scholar.[17] It was received favourably by the national heraldic commission[18]
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The blazon given on the region's website is Écartelé : au premier d’or au gonfanon de gueules bordé de sinople (Auvergne) ; au deuxième de gueules à la croix d’argent (Savoie) ; au troisième de gueules au lion d’argent (Lyonnais) ; au quatrième d’or au dauphin d’azur, crêté, barbé, loré, peautré et oreillé de gueules (Dauphiné)., which translates roughly to: Quarterly: 1st Or a gonfalon Gules bordered Vert (Auverne); 2nd Gules a cross argent (Savoie); 3rd Gules a lion argent (Lyonnais); 4th Or a dolphin Azure, crested, bearded, lored, and oreilled Gules.[17]
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Official coat of arms, released in 2017 and officialized in 2018.
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First flag, which appeared at the same time as the coat of arms.
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Flag officialized in 2018
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In Francoprovençal and in Occitan, two of the three languages that are historically spoken in the region, the name is pronounced:[Note 2]:
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The Auvergne-Rhône-Alpes administrative region covers an area of 69 711 km2 in the centre and east of the south of France. It is a collection of regions of diverse topographies, climates, natural resources, cultures, folklore, architecture, and languages. It is bordered by five other administrative regions: Bourgogne-Franche-Comté to the north, Centre-Val de Loire to the northeast, Nouvelle-Aquitaine to the west, Occitanie to the south-west, and Provence-Alpes-Côte d'Azur to the south-east. It is also bordered by Italy (Aosta Valley and Piedmont) to the east and Switzerland (Cantons of Geneva, Valais, and Vaud) to the north-east.
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Extreme points:
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Auvergne-Rhône-Alpes comprises twelve departments: Ain, Allier, Ardèche, Cantal, Drôme, Haute-Loire, Haute-Savoie, Isère, Loire, Puy-de-Dôme, Rhône, Savoie
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The Gross domestic product (GDP) of the region was 270.0 billion euros in 2018, accounting for 11.5% of French economic output. GDP per capita adjusted for purchasing power was 30,200 euros or 100% of the EU27 average in the same year. The GDP per employee was 109% of the EU average.[19]
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The region is governed by the regional council of Auvergne-Rhône-Alpes consisting of 204 members. The current regional council was elected in regional elections on 6 and 13 December 2015, with the list of Laurent Wauquiez consisting of The Republicans (LR), the Democratic Movement (MoDem), and Union of Democrats and Independents (UDI) securing an absolute majority of 113 seats.[20][21]
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Proteins are large biomolecules, or macromolecules, consisting of one or more long chains of amino acid residues. Proteins perform a vast array of functions within organisms, including catalysing metabolic reactions, DNA replication, responding to stimuli, providing structure to cells, and organisms, and transporting molecules from one location to another. Proteins differ from one another primarily in their sequence of amino acids, which is dictated by the nucleotide sequence of their genes, and which usually results in protein folding into a specific 3D structure that determines its activity.
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A linear chain of amino acid residues is called a polypeptide. A protein contains at least one long polypeptide. Short polypeptides, containing less than 20–30 residues, are rarely considered to be proteins and are commonly called peptides, or sometimes oligopeptides. The individual amino acid residues are bonded together by peptide bonds and adjacent amino acid residues. The sequence of amino acid residues in a protein is defined by the sequence of a gene, which is encoded in the genetic code. In general, the genetic code specifies 20 standard amino acids; but in certain organisms the genetic code can include selenocysteine and—in certain archaea—pyrrolysine. Shortly after or even during synthesis, the residues in a protein are often chemically modified by post-translational modification, which alters the physical and chemical properties, folding, stability, activity, and ultimately, the function of the proteins. Some proteins have non-peptide groups attached, which can be called prosthetic groups or cofactors. Proteins can also work together to achieve a particular function, and they often associate to form stable protein complexes.
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Once formed, proteins only exist for a certain period and are then degraded and recycled by the cell's machinery through the process of protein turnover. A protein's lifespan is measured in terms of its half-life and covers a wide range. They can exist for minutes or years with an average lifespan of 1–2 days in mammalian cells. Abnormal or misfolded proteins are degraded more rapidly either due to being targeted for destruction or due to being unstable.
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Like other biological macromolecules such as polysaccharides and nucleic acids, proteins are essential parts of organisms and participate in virtually every process within cells. Many proteins are enzymes that catalyse biochemical reactions and are vital to metabolism. Proteins also have structural or mechanical functions, such as actin and myosin in muscle and the proteins in the cytoskeleton, which form a system of scaffolding that maintains cell shape. Other proteins are important in cell signaling, immune responses, cell adhesion, and the cell cycle. In animals, proteins are needed in the diet to provide the essential amino acids that cannot be synthesized. Digestion breaks the proteins down for use in the metabolism.
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Proteins may be purified from other cellular components using a variety of techniques such as ultracentrifugation, precipitation, electrophoresis, and chromatography; the advent of genetic engineering has made possible a number of methods to facilitate purification. Methods commonly used to study protein structure and function include immunohistochemistry, site-directed mutagenesis, X-ray crystallography, nuclear magnetic resonance and mass spectrometry.
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Proteins were recognized as a distinct class of biological molecules in the eighteenth century by Antoine Fourcroy and others, distinguished by the molecules' ability to coagulate or flocculate under treatments with heat or acid.[1] Noted examples at the time included albumin from egg whites, blood serum albumin, fibrin, and wheat gluten.
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Proteins were first described by the Dutch chemist Gerardus Johannes Mulder and named by the Swedish chemist Jöns Jacob Berzelius in 1838.[2][3] Mulder carried out elemental analysis of common proteins and found that nearly all proteins had the same empirical formula, C400H620N100O120P1S1.[4] He came to the erroneous conclusion that they might be composed of a single type of (very large) molecule. The term "protein" to describe these molecules was proposed by Mulder's associate Berzelius; protein is derived from the Greek word πρώτειος (proteios), meaning "primary",[5] "in the lead", or "standing in front",[6] + -in. Mulder went on to identify the products of protein degradation such as the amino acid leucine for which he found a (nearly correct) molecular weight of 131 Da.[4] Prior to "protein", other names were used, like "albumins" or "albuminous materials" (Eiweisskörper, in German).[7]
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Early nutritional scientists such as the German Carl von Voit believed that protein was the most important nutrient for maintaining the structure of the body, because it was generally believed that "flesh makes flesh."[8] Karl Heinrich Ritthausen extended known protein forms with the identification of glutamic acid. At the Connecticut Agricultural Experiment Station a detailed review of the vegetable proteins was compiled by Thomas Burr Osborne. Working with Lafayette Mendel and applying Liebig's law of the minimum in feeding laboratory rats, the nutritionally essential amino acids were established. The work was continued and communicated by William Cumming Rose. The understanding of proteins as polypeptides came through the work of Franz Hofmeister and Hermann Emil Fischer in 1902.[9][10] The central role of proteins as enzymes in living organisms was not fully appreciated until 1926, when James B. Sumner showed that the enzyme urease was in fact a protein.[11]
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The difficulty in purifying proteins in large quantities made them very difficult for early protein biochemists to study. Hence, early studies focused on proteins that could be purified in large quantities, e.g., those of blood, egg white, various toxins, and digestive/metabolic enzymes obtained from slaughterhouses. In the 1950s, the Armour Hot Dog Co. purified 1 kg of pure bovine pancreatic ribonuclease A and made it freely available to scientists; this gesture helped ribonuclease A become a major target for biochemical study for the following decades.[4]
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Linus Pauling is credited with the successful prediction of regular protein secondary structures based on hydrogen bonding, an idea first put forth by William Astbury in 1933.[12] Later work by Walter Kauzmann on denaturation,[13][14] based partly on previous studies by Kaj Linderstrøm-Lang,[15] contributed an understanding of protein folding and structure mediated by hydrophobic interactions.
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The first protein to be sequenced was insulin, by Frederick Sanger, in 1949. Sanger correctly determined the amino acid sequence of insulin, thus conclusively demonstrating that proteins consisted of linear polymers of amino acids rather than branched chains, colloids, or cyclols.[16] He won the Nobel Prize for this achievement in 1958.[17]
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The first protein structures to be solved were hemoglobin and myoglobin, by Max Perutz and Sir John Cowdery Kendrew, respectively, in 1958.[18][19] As of 2017[update], the Protein Data Bank has over 126,060 atomic-resolution structures of proteins.[20] In more recent times, cryo-electron microscopy of large macromolecular assemblies[21] and computational protein structure prediction of small protein domains[22] are two methods approaching atomic resolution.
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Most proteins consist of linear polymers built from series of up to 20 different L-α- amino acids. All proteinogenic amino acids possess common structural features, including an α-carbon to which an amino group, a carboxyl group, and a variable side chain are bonded. Only proline differs from this basic structure as it contains an unusual ring to the N-end amine group, which forces the CO–NH amide moiety into a fixed conformation.[23] The side chains of the standard amino acids, detailed in the list of standard amino acids, have a great variety of chemical structures and properties; it is the combined effect of all of the amino acid side chains in a protein that ultimately determines its three-dimensional structure and its chemical reactivity.[24]
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The amino acids in a polypeptide chain are linked by peptide bonds. Once linked in the protein chain, an individual amino acid is called a residue, and the linked series of carbon, nitrogen, and oxygen atoms are known as the main chain or protein backbone.[25]
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The peptide bond has two resonance forms that contribute some double-bond character and inhibit rotation around its axis, so that the alpha carbons are roughly coplanar. The other two dihedral angles in the peptide bond determine the local shape assumed by the protein backbone.[26] The end with a free amino group is known as the N-terminus or amino terminus, whereas the end of the protein with a free carboxyl group is known as the C-terminus or carboxy terminus (the sequence of the protein is written from N-terminus to C-terminus, from left to right).
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The words protein, polypeptide, and peptide are a little ambiguous and can overlap in meaning. Protein is generally used to refer to the complete biological molecule in a stable conformation, whereas peptide is generally reserved for a short amino acid oligomers often lacking a stable 3D structure. But the boundary between the two is not well defined and usually lies near 20–30 residues.[27] Polypeptide can refer to any single linear chain of amino acids, usually regardless of length, but often implies an absence of a defined conformation.
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Proteins can interact with many types of molecules, including with other proteins, with lipids, with carboyhydrates, and with DNA.[28][29][30][31]
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It has been estimated that average-sized bacteria contain about 2 million proteins per cell (e.g. E. coli and Staphylococcus aureus). Smaller bacteria, such as Mycoplasma or spirochetes contain fewer molecules, on the order of 50,000 to 1 million. By contrast, eukaryotic cells are larger and thus contain much more protein. For instance, yeast cells have been estimated to contain about 50 million proteins and human cells on the order of 1 to 3 billion.[32] The concentration of individual protein copies ranges from a few molecules per cell up to 20 million.[33] Not all genes coding proteins are expressed in most cells and their number depends on, for example, cell type and external stimuli. For instance, of the 20,000 or so proteins encoded by the human genome, only 6,000 are detected in lymphoblastoid cells.[34] Moreover, the number of proteins the genome encodes correlates well with the organism complexity. Eukaryotes have 15,000, bacteria have 3,200, archaea have 2,400, and viruses have 42 proteins on average coded in their respective genomes.[35]
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Proteins are assembled from amino acids using information encoded in genes. Each protein has its own unique amino acid sequence that is specified by the nucleotide sequence of the gene encoding this protein. The genetic code is a set of three-nucleotide sets called codons and each three-nucleotide combination designates an amino acid, for example AUG (adenine–uracil–guanine) is the code for methionine. Because DNA contains four nucleotides, the total number of possible codons is 64; hence, there is some redundancy in the genetic code, with some amino acids specified by more than one codon.[36] Genes encoded in DNA are first transcribed into pre-messenger RNA (mRNA) by proteins such as RNA polymerase. Most organisms then process the pre-mRNA (also known as a primary transcript) using various forms of Post-transcriptional modification to form the mature mRNA, which is then used as a template for protein synthesis by the ribosome. In prokaryotes the mRNA may either be used as soon as it is produced, or be bound by a ribosome after having moved away from the nucleoid. In contrast, eukaryotes make mRNA in the cell nucleus and then translocate it across the nuclear membrane into the cytoplasm, where protein synthesis then takes place. The rate of protein synthesis is higher in prokaryotes than eukaryotes and can reach up to 20 amino acids per second.[37]
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The process of synthesizing a protein from an mRNA template is known as translation. The mRNA is loaded onto the ribosome and is read three nucleotides at a time by matching each codon to its base pairing anticodon located on a transfer RNA molecule, which carries the amino acid corresponding to the codon it recognizes. The enzyme aminoacyl tRNA synthetase "charges" the tRNA molecules with the correct amino acids. The growing polypeptide is often termed the nascent chain. Proteins are always biosynthesized from N-terminus to C-terminus.[36]
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The size of a synthesized protein can be measured by the number of amino acids it contains and by its total molecular mass, which is normally reported in units of daltons (synonymous with atomic mass units), or the derivative unit kilodalton (kDa). The average size of a protein increases from Archaea to Bacteria to Eukaryote (283, 311, 438 residues and 31, 34, 49 kDa respectively) due to a bigger number of protein domains constituting proteins in higher organisms.[35] For instance, yeast proteins are on average 466 amino acids long and 53 kDa in mass.[27] The largest known proteins are the titins, a component of the muscle sarcomere, with a molecular mass of almost 3,000 kDa and a total length of almost 27,000 amino acids.[38]
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Short proteins can also be synthesized chemically by a family of methods known as peptide synthesis, which rely on organic synthesis techniques such as chemical ligation to produce peptides in high yield.[39] Chemical synthesis allows for the introduction of non-natural amino acids into polypeptide chains, such as attachment of fluorescent probes to amino acid side chains.[40] These methods are useful in laboratory biochemistry and cell biology, though generally not for commercial applications. Chemical synthesis is inefficient for polypeptides longer than about 300 amino acids, and the synthesized proteins may not readily assume their native tertiary structure. Most chemical synthesis methods proceed from C-terminus to N-terminus, opposite the biological reaction.[41]
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Most proteins fold into unique 3D structures. The shape into which a protein naturally folds is known as its native conformation.[42] Although many proteins can fold unassisted, simply through the chemical properties of their amino acids, others require the aid of molecular chaperones to fold into their native states.[43] Biochemists often refer to four distinct aspects of a protein's structure:[44]
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Proteins are not entirely rigid molecules. In addition to these levels of structure, proteins may shift between several related structures while they perform their functions. In the context of these functional rearrangements, these tertiary or quaternary structures are usually referred to as "conformations", and transitions between them are called conformational changes. Such changes are often induced by the binding of a substrate molecule to an enzyme's active site, or the physical region of the protein that participates in chemical catalysis. In solution proteins also undergo variation in structure through thermal vibration and the collision with other molecules.[45]
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Proteins can be informally divided into three main classes, which correlate with typical tertiary structures: globular proteins, fibrous proteins, and membrane proteins. Almost all globular proteins are soluble and many are enzymes. Fibrous proteins are often structural, such as collagen, the major component of connective tissue, or keratin, the protein component of hair and nails. Membrane proteins often serve as receptors or provide channels for polar or charged molecules to pass through the cell membrane.[46]
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A special case of intramolecular hydrogen bonds within proteins, poorly shielded from water attack and hence promoting their own dehydration, are called dehydrons.[47]
|
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Many proteins are composed of several protein domains, i.e. segments of a protein that fold into distinct structural units. Domains usually also have specific functions, such as enzymatic activities (e.g. kinase) or they serve as binding modules (e.g. the SH3 domain binds to proline-rich sequences in other proteins).
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Short amino acid sequences within proteins often act as recognition sites for other proteins.[48] For instance, SH3 domains typically bind to short PxxP motifs (i.e. 2 prolines [P], separated by two unspecified amino acids [x], although the surrounding amino acids may determine the exact binding specificity). Many such motifs has been collected in the Eukaryotic Linear Motif (ELM) database.
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Proteins are the chief actors within the cell, said to be carrying out the duties specified by the information encoded in genes.[27] With the exception of certain types of RNA, most other biological molecules are relatively inert elements upon which proteins act. Proteins make up half the dry weight of an Escherichia coli cell, whereas other macromolecules such as DNA and RNA make up only 3% and 20%, respectively.[49] The set of proteins expressed in a particular cell or cell type is known as its proteome.
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The chief characteristic of proteins that also allows their diverse set of functions is their ability to bind other molecules specifically and tightly. The region of the protein responsible for binding another molecule is known as the binding site and is often a depression or "pocket" on the molecular surface. This binding ability is mediated by the tertiary structure of the protein, which defines the binding site pocket, and by the chemical properties of the surrounding amino acids' side chains. Protein binding can be extraordinarily tight and specific; for example, the ribonuclease inhibitor protein binds to human angiogenin with a sub-femtomolar dissociation constant (<10−15 M) but does not bind at all to its amphibian homolog onconase (>1 M). Extremely minor chemical changes such as the addition of a single methyl group to a binding partner can sometimes suffice to nearly eliminate binding; for example, the aminoacyl tRNA synthetase specific to the amino acid valine discriminates against the very similar side chain of the amino acid isoleucine.[50]
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Proteins can bind to other proteins as well as to small-molecule substrates. When proteins bind specifically to other copies of the same molecule, they can oligomerize to form fibrils; this process occurs often in structural proteins that consist of globular monomers that self-associate to form rigid fibers. Protein–protein interactions also regulate enzymatic activity, control progression through the cell cycle, and allow the assembly of large protein complexes that carry out many closely related reactions with a common biological function. Proteins can also bind to, or even be integrated into, cell membranes. The ability of binding partners to induce conformational changes in proteins allows the construction of enormously complex signaling networks.[51]
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As interactions between proteins are reversible, and depend heavily on the availability of different groups of partner proteins to form aggregates that are capable to carry out discrete sets of function, study of the interactions between specific proteins is a key to understand important aspects of cellular function, and ultimately the properties that distinguish particular cell types.[52][53]
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The best-known role of proteins in the cell is as enzymes, which catalyse chemical reactions. Enzymes are usually highly specific and accelerate only one or a few chemical reactions. Enzymes carry out most of the reactions involved in metabolism, as well as manipulating DNA in processes such as DNA replication, DNA repair, and transcription. Some enzymes act on other proteins to add or remove chemical groups in a process known as posttranslational modification. About 4,000 reactions are known to be catalysed by enzymes.[54] The rate acceleration conferred by enzymatic catalysis is often enormous—as much as 1017-fold increase in rate over the uncatalysed reaction in the case of orotate decarboxylase (78 million years without the enzyme, 18 milliseconds with the enzyme).[55]
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The molecules bound and acted upon by enzymes are called substrates. Although enzymes can consist of hundreds of amino acids, it is usually only a small fraction of the residues that come in contact with the substrate, and an even smaller fraction—three to four residues on average—that are directly involved in catalysis.[56] The region of the enzyme that binds the substrate and contains the catalytic residues is known as the active site.
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Dirigent proteins are members of a class of proteins that dictate the stereochemistry of a compound synthesized by other enzymes.[57]
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Many proteins are involved in the process of cell signaling and signal transduction. Some proteins, such as insulin, are extracellular proteins that transmit a signal from the cell in which they were synthesized to other cells in distant tissues. Others are membrane proteins that act as receptors whose main function is to bind a signaling molecule and induce a biochemical response in the cell. Many receptors have a binding site exposed on the cell surface and an effector domain within the cell, which may have enzymatic activity or may undergo a conformational change detected by other proteins within the cell.[58]
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Antibodies are protein components of an adaptive immune system whose main function is to bind antigens, or foreign substances in the body, and target them for destruction. Antibodies can be secreted into the extracellular environment or anchored in the membranes of specialized B cells known as plasma cells. Whereas enzymes are limited in their binding affinity for their substrates by the necessity of conducting their reaction, antibodies have no such constraints. An antibody's binding affinity to its target is extraordinarily high.[59]
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Many ligand transport proteins bind particular small biomolecules and transport them to other locations in the body of a multicellular organism. These proteins must have a high binding affinity when their ligand is present in high concentrations, but must also release the ligand when it is present at low concentrations in the target tissues. The canonical example of a ligand-binding protein is haemoglobin, which transports oxygen from the lungs to other organs and tissues in all vertebrates and has close homologs in every biological kingdom.[60] Lectins are sugar-binding proteins which are highly specific for their sugar moieties. Lectins typically play a role in biological recognition phenomena involving cells and proteins.[61] Receptors and hormones are highly specific binding proteins.
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Transmembrane proteins can also serve as ligand transport proteins that alter the permeability of the cell membrane to small molecules and ions. The membrane alone has a hydrophobic core through which polar or charged molecules cannot diffuse. Membrane proteins contain internal channels that allow such molecules to enter and exit the cell. Many ion channel proteins are specialized to select for only a particular ion; for example, potassium and sodium channels often discriminate for only one of the two ions.[62]
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Structural proteins confer stiffness and rigidity to otherwise-fluid biological components. Most structural proteins are fibrous proteins; for example, collagen and elastin are critical components of connective tissue such as cartilage, and keratin is found in hard or filamentous structures such as hair, nails, feathers, hooves, and some animal shells.[63] Some globular proteins can also play structural functions, for example, actin and tubulin are globular and soluble as monomers, but polymerize to form long, stiff fibers that make up the cytoskeleton, which allows the cell to maintain its shape and size.
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Other proteins that serve structural functions are motor proteins such as myosin, kinesin, and dynein, which are capable of generating mechanical forces. These proteins are crucial for cellular motility of single celled organisms and the sperm of many multicellular organisms which reproduce sexually. They also generate the forces exerted by contracting muscles[64] and play essential roles in intracellular transport.
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The activities and structures of proteins may be examined in vitro, in vivo, and in silico. In vitro studies of purified proteins in controlled environments are useful for learning how a protein carries out its function: for example, enzyme kinetics studies explore the chemical mechanism of an enzyme's catalytic activity and its relative affinity for various possible substrate molecules. By contrast, in vivo experiments can provide information about the physiological role of a protein in the context of a cell or even a whole organism. In silico studies use computational methods to study proteins.
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To perform in vitro analysis, a protein must be purified away from other cellular components. This process usually begins with cell lysis, in which a cell's membrane is disrupted and its internal contents released into a solution known as a crude lysate. The resulting mixture can be purified using ultracentrifugation, which fractionates the various cellular components into fractions containing soluble proteins; membrane lipids and proteins; cellular organelles, and nucleic acids. Precipitation by a method known as salting out can concentrate the proteins from this lysate. Various types of chromatography are then used to isolate the protein or proteins of interest based on properties such as molecular weight, net charge and binding affinity.[65] The level of purification can be monitored using various types of gel electrophoresis if the desired protein's molecular weight and isoelectric point are known, by spectroscopy if the protein has distinguishable spectroscopic features, or by enzyme assays if the protein has enzymatic activity. Additionally, proteins can be isolated according to their charge using electrofocusing.[66]
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For natural proteins, a series of purification steps may be necessary to obtain protein sufficiently pure for laboratory applications. To simplify this process, genetic engineering is often used to add chemical features to proteins that make them easier to purify without affecting their structure or activity. Here, a "tag" consisting of a specific amino acid sequence, often a series of histidine residues (a "His-tag"), is attached to one terminus of the protein. As a result, when the lysate is passed over a chromatography column containing nickel, the histidine residues ligate the nickel and attach to the column while the untagged components of the lysate pass unimpeded. A number of different tags have been developed to help researchers purify specific proteins from complex mixtures.[67]
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The study of proteins in vivo is often concerned with the synthesis and localization of the protein within the cell. Although many intracellular proteins are synthesized in the cytoplasm and membrane-bound or secreted proteins in the endoplasmic reticulum, the specifics of how proteins are targeted to specific organelles or cellular structures is often unclear. A useful technique for assessing cellular localization uses genetic engineering to express in a cell a fusion protein or chimera consisting of the natural protein of interest linked to a "reporter" such as green fluorescent protein (GFP).[68] The fused protein's position within the cell can be cleanly and efficiently visualized using microscopy,[69] as shown in the figure opposite.
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Other methods for elucidating the cellular location of proteins requires the use of known compartmental markers for regions such as the ER, the Golgi, lysosomes or vacuoles, mitochondria, chloroplasts, plasma membrane, etc. With the use of fluorescently tagged versions of these markers or of antibodies to known markers, it becomes much simpler to identify the localization of a protein of interest. For example, indirect immunofluorescence will allow for fluorescence colocalization and demonstration of location. Fluorescent dyes are used to label cellular compartments for a similar purpose.[70]
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Other possibilities exist, as well. For example, immunohistochemistry usually utilizes an antibody to one or more proteins of interest that are conjugated to enzymes yielding either luminescent or chromogenic signals that can be compared between samples, allowing for localization information. Another applicable technique is cofractionation in sucrose (or other material) gradients using isopycnic centrifugation.[71] While this technique does not prove colocalization of a compartment of known density and the protein of interest, it does increase the likelihood, and is more amenable to large-scale studies.
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Finally, the gold-standard method of cellular localization is immunoelectron microscopy. This technique also uses an antibody to the protein of interest, along with classical electron microscopy techniques. The sample is prepared for normal electron microscopic examination, and then treated with an antibody to the protein of interest that is conjugated to an extremely electro-dense material, usually gold. This allows for the localization of both ultrastructural details as well as the protein of interest.[72]
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Through another genetic engineering application known as site-directed mutagenesis, researchers can alter the protein sequence and hence its structure, cellular localization, and susceptibility to regulation. This technique even allows the incorporation of unnatural amino acids into proteins, using modified tRNAs,[73] and may allow the rational design of new proteins with novel properties.[74]
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The total complement of proteins present at a time in a cell or cell type is known as its proteome, and the study of such large-scale data sets defines the field of proteomics, named by analogy to the related field of genomics. Key experimental techniques in proteomics include 2D electrophoresis,[75] which allows the separation of many proteins, mass spectrometry,[76] which allows rapid high-throughput identification of proteins and sequencing of peptides (most often after in-gel digestion), protein microarrays, which allow the detection of the relative levels of the various proteins present in a cell, and two-hybrid screening, which allows the systematic exploration of protein–protein interactions.[77] The total complement of biologically possible such interactions is known as the interactome.[78] A systematic attempt to determine the structures of proteins representing every possible fold is known as structural genomics.[79]
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A vast array of computational methods have been developed to analyze the structure, function and evolution of proteins.
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The development of such tools has been driven by the large amount of genomic and proteomic data available for a variety of organisms, including the human genome. It is simply impossible to study all proteins experimentally, hence only a few are subjected to laboratory experiments while computational tools are used to extrapolate to similar proteins. Such homologous proteins can be efficiently identified in distantly related organisms by sequence alignment. Genome and gene sequences can be searched by a variety of tools for certain properties. Sequence profiling tools can find restriction enzyme sites, open reading frames in nucleotide sequences, and predict secondary structures. Phylogenetic trees can be constructed and evolutionary hypotheses developed using special software like ClustalW regarding the ancestry of modern organisms and the genes they express. The field of bioinformatics is now indispensable for the analysis of genes and proteins.
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Discovering the tertiary structure of a protein, or the quaternary structure of its complexes, can provide important clues about how the protein performs its function and how it can be affected, i.e. in drug design. As proteins are too small to be seen under a light microscope, other methods have to be employed to determine their structure. Common experimental methods include X-ray crystallography and NMR spectroscopy, both of which can produce structural information at atomic resolution. However, NMR experiments are able to provide information from which a subset of distances between pairs of atoms can be estimated, and the final possible conformations for a protein are determined by solving a distance geometry problem. Dual polarisation interferometry is a quantitative analytical method for measuring the overall protein conformation and conformational changes due to interactions or other stimulus. Circular dichroism is another laboratory technique for determining internal β-sheet / α-helical composition of proteins. Cryoelectron microscopy is used to produce lower-resolution structural information about very large protein complexes, including assembled viruses;[80] a variant known as electron crystallography can also produce high-resolution information in some cases, especially for two-dimensional crystals of membrane proteins.[81] Solved structures are usually deposited in the Protein Data Bank (PDB), a freely available resource from which structural data about thousands of proteins can be obtained in the form of Cartesian coordinates for each atom in the protein.[82]
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Many more gene sequences are known than protein structures. Further, the set of solved structures is biased toward proteins that can be easily subjected to the conditions required in X-ray crystallography, one of the major structure determination methods. In particular, globular proteins are comparatively easy to crystallize in preparation for X-ray crystallography. Membrane proteins and large protein complexes, by contrast, are difficult to crystallize and are underrepresented in the PDB.[83] Structural genomics initiatives have attempted to remedy these deficiencies by systematically solving representative structures of major fold classes. Protein structure prediction methods attempt to provide a means of generating a plausible structure for proteins whose structures have not been experimentally determined.[84]
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Complementary to the field of structural genomics, protein structure prediction develops efficient mathematical models of proteins to computationally predict the molecular formations in theory, instead of detecting structures with laboratory observation.[85] The most successful type of structure prediction, known as homology modeling, relies on the existence of a "template" structure with sequence similarity to the protein being modeled; structural genomics' goal is to provide sufficient representation in solved structures to model most of those that remain.[86] Although producing accurate models remains a challenge when only distantly related template structures are available, it has been suggested that sequence alignment is the bottleneck in this process, as quite accurate models can be produced if a "perfect" sequence alignment is known.[87] Many structure prediction methods have served to inform the emerging field of protein engineering, in which novel protein folds have already been designed.[88] A more complex computational problem is the prediction of intermolecular interactions, such as in molecular docking and protein–protein interaction prediction.[89]
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Mathematical models to simulate dynamic processes of protein folding and binding involve molecular mechanics, in particular, molecular dynamics. Monte Carlo techniques facilitate the computations, which exploit advances in parallel and distributed computing (for example, the Folding@home project[90] which performs molecular modeling on GPUs). In silico simulations discovered the folding of small α-helical protein domains such as the villin headpiece[91] and the HIV accessory protein.[92] Hybrid methods combining standard molecular dynamics with quantum mechanical mathematics explored the electronic states of rhodopsins.[93]
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Many proteins (in Eucaryota ~33%) contain large unstructured but biologically functional segments and can be classified as intrinsically disordered proteins.[94] Predicting and analysing protein disorder is, therefore, an important part of protein structure characterisation.[95]
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Most microorganisms and plants can biosynthesize all 20 standard amino acids, while animals (including humans) must obtain some of the amino acids from the diet.[49] The amino acids that an organism cannot synthesize on its own are referred to as essential amino acids. Key enzymes that synthesize certain amino acids are not present in animals—such as aspartokinase, which catalyses the first step in the synthesis of lysine, methionine, and threonine from aspartate. If amino acids are present in the environment, microorganisms can conserve energy by taking up the amino acids from their surroundings and downregulating their biosynthetic pathways.
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In animals, amino acids are obtained through the consumption of foods containing protein. Ingested proteins are then broken down into amino acids through digestion, which typically involves denaturation of the protein through exposure to acid and hydrolysis by enzymes called proteases. Some ingested amino acids are used for protein biosynthesis, while others are converted to glucose through gluconeogenesis, or fed into the citric acid cycle. This use of protein as a fuel is particularly important under starvation conditions as it allows the body's own proteins to be used to support life, particularly those found in muscle.[96]
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In animals such as dogs and cats, protein maintains the health and quality of the skin by promoting hair follicle growth and keratinization, and thus reducing the likelihood of skin problems producing malodours.[97] Poor-quality proteins also have a role regarding gastrointestinal health, increasing the potential for flatulence and odorous compounds in dogs because when proteins reach the colon in an undigested state, they are fermented producing hydrogen sulfide gas, indole, and skatole.[98] Dogs and cats digest animal proteins better than those from plants, but products of low-quality animal origin are poorly digested, including skin, feathers, and connective tissue.[98]
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Supergroups[1] and typical phyla
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Many others;classification varies
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A protist (/ˈproʊtɪst/) is any eukaryotic organism (one with cells containing a nucleus) that is not an animal, plant, or fungus. While it is likely that protists share a common ancestor (the last eukaryotic common ancestor),[2] the exclusion of other eukaryotes means that protists do not form a natural group, or clade.[a] So some protists may be more closely related to animals, plants, or fungi than they are to other protists; however, like algae, invertebrates, or protozoans, the grouping is used for convenience. The study of protists is termed protistology.[3]
|
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+
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+
The classification of a separate kingdom to the animals and plants was first proposed by John Hogg in 1860 as the kingdom Protoctista; in 1866 Ernst Haeckel also proposed a third kingdom Protista as "the kingdom of primitive forms".[4] Originally these also included prokaryotes, but with time these would be removed to a fourth kingdom Monera.[b] In the popular five-kingdom scheme proposed by Robert Whittaker in 1969, the Protista was defined as eukaryotic "organisms which are unicellular or unicellular-colonial and which form no tissues", and the fifth kingdom Fungi was established.[5][6][c] In the five-kingdom system of Lynn Margulis, the term protist is reserved for microscopic organisms, while the more inclusive kingdom Protoctista (or protoctists) included certain large multicellular eukaryotes, such as kelp, red algae and slime molds.[9] Others use the term protist interchangeably with Margulis's protoctist, to encompass both single-celled and multicellular eukaryotes, including those that form specialized tissues but do not fit into any of the other traditional kingdoms.[10]
|
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+
|
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+
Besides their relatively simple levels of organization, protists do not necessarily have much in common.[11] When used, the term "protists" is now considered to mean a paraphyletic assemblage of similar-appearing but diverse taxa (biological groups); these taxa do not have an exclusive common ancestor beyond being composed of eukaryotes, and have different life cycles, trophic levels, modes of locomotion and cellular structures.[12][13] Examples of protists include:[14] amoebas (including nucleariids and Foraminifera); choanaflagellates; ciliates; diatoms; dinoflagellates; Giardia; Plasmodium (which causes malaria); oomycetes (including Phytophthora, the cause of the Great Famine of Ireland); and slime molds. These examples are unicellular, although oomycetes can form filaments, and slime molds can aggregate.
|
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+
|
11 |
+
In cladistic systems (classifications based on common ancestry), there are no equivalents to the taxa Protista or Protoctista, as both terms refer to a paraphyletic group that spans the entire eukaryotic tree of life. In cladistic classification, the contents of Protista are mostly distributed among various supergroups: examples include the SAR supergroup (of stramenopiles or heterokonts, alveolates, and Rhizaria); Archaeplastida (or Plantae sensu lato); Excavata (which is mostly unicellular flagellates); and Opisthokonta (which commonly includes unicellular flagellates, but also animals and fungi). "Protista", "Protoctista", and "Protozoa" are therefore considered obsolete. However, the term "protist" continues to be used informally as a catch-all term for eukayotic organisms that aren't within other traditional kingdoms. For example, the word "protist pathogen" may be used to denote any disease-causing organism that is not plant, animal, fungal, prokaryotic, viral, or subviral.[15]
|
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+
|
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+
The term protista was first used by Ernst Haeckel in 1866. Protists were traditionally subdivided into several groups based on similarities to the "higher" kingdoms such as:[4]
|
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+
|
15 |
+
Some protists, sometimes called ambiregnal protists, have been considered to be both protozoa and algae or fungi (e.g., slime molds and flagellated algae), and names for these have been published under either or both of the ICN and the ICZN.[16][17] Conflicts, such as these – for example the dual-classification of Euglenids and Dinobryons, which are mixotrophic – is an example of why the kingdom Protista was adopted.
|
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+
|
17 |
+
These traditional subdivisions, largely based on superficial commonalities, have been replaced by classifications based on phylogenetics (evolutionary relatedness among organisms).
|
18 |
+
Molecular analyses in modern taxonomy have been used to redistribute former members of this group into diverse and sometimes distantly related phyla. For instance, the water molds are now considered to be closely related to photosynthetic organisms such as Brown algae and Diatoms, the slime molds are grouped mainly under Amoebozoa, and the Amoebozoa itself includes only a subset of the "Amoeba" group, and significant number of erstwhile "Amoeboid" genera are distributed among Rhizarians and other Phyla.
|
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+
|
20 |
+
However, the older terms are still used as informal names to describe the morphology and ecology of various protists. For example, the term protozoa is used to refer to heterotrophic species of protists that do not form filaments.
|
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+
|
22 |
+
Among the pioneers in the study of the protists, which were almost ignored by Linnaeus except for some genera (e.g., Vorticella, Chaos, Volvox, Corallina, Conferva, Ulva, Chara, Fucus)[18][19] were Leeuwenhoek, O. F. Müller, C. G. Ehrenberg and Félix Dujardin.[20] The first groups used to classify microscopic organism were the Animalcules and the Infusoria.[21] In 1818, the German naturalist Georg August Goldfuss introduced the word Protozoa to refer to organisms such as ciliates and corals.[22][4] After the cell theory of Schwann and Schleiden (1838–39), this group was modified in 1848 by Carl von Siebold to include only animal-like unicellular organisms, such as foraminifera and amoebae.[23] The formal taxonomic category Protoctista was first proposed in the early 1860s by John Hogg, who argued that the protists should include what he saw as primitive unicellular forms of both plants and animals. He defined the Protoctista as a "fourth kingdom of nature", in addition to the then-traditional kingdoms of plants, animals and minerals.[24][4] The kingdom of minerals was later removed from taxonomy in 1866 by Ernst Haeckel, leaving plants, animals, and the protists (Protista), defined as a "kingdom of primitive forms".[25][26]
|
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+
|
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+
In 1938, Herbert Copeland resurrected Hogg's label, arguing that Haeckel's term Protista included anucleated microbes such as bacteria, which the term "Protoctista" (literally meaning "first established beings") did not. In contrast, Copeland's term included nucleated eukaryotes such as diatoms, green algae and fungi.[27] This classification was the basis for Whittaker's later definition of Fungi, Animalia, Plantae and Protista as the four kingdoms of life.[7] The kingdom Protista was later modified to separate prokaryotes into the separate kingdom of Monera, leaving the protists as a group of eukaryotic microorganisms.[5] These five kingdoms remained the accepted classification until the development of molecular phylogenetics in the late 20th century, when it became apparent that neither protists nor monera were single groups of related organisms (they were not monophyletic groups).[28]
|
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+
|
26 |
+
Systematists today do not treat Protista as a formal taxon, but the term "protist" is still commonly used for convenience in two ways.[29] The most popular contemporary definition is a phylogenetic one, that identifies a paraphyletic group:[30] a protist is any eukaryote that is not an animal, (land) plant, or (true) fungus; this definition[31] excludes many unicellular groups, like the Microsporidia (fungi), many Chytridiomycetes (fungi), and yeasts (fungi), and also a non-unicellular group included in Protista in the past, the Myxozoa (animal).[32] Some systematists[who?] judge paraphyletic taxa acceptable, and use Protista in this sense as a formal taxon (as found in some secondary textbooks, for pedagogical purpose).[citation needed]
|
27 |
+
|
28 |
+
The other definition describes protists primarily by functional or biological criteria: protists are essentially those eukaryotes that are never multicellular,[29] that either exist as independent cells, or if they occur in colonies, do not show differentiation into tissues (but vegetative cell differentiation may occur restricted to sexual reproduction, alternate vegetative morphology, and quiescent or resistant stages, such as cysts);[33] this definition excludes many brown, multicellular red and green algae, which may have tissues.
|
29 |
+
|
30 |
+
The taxonomy of protists is still changing. Newer classifications attempt to present monophyletic groups based on morphological (especially ultrastructural),[34][35][36] biochemical (chemotaxonomy)[37][38] and DNA sequence (molecular research) information.[39][40] However, there are sometimes discordances between molecular and morphological investigations; these can be categorized as two types: (i) one morphology, multiple lineages (e.g. morphological convergence, cryptic species) and (ii) one lineage, multiple morphologies (e.g. phenotypic plasticity, multiple life-cycle stages).[41]
|
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+
|
32 |
+
Because the protists as a whole are paraphyletic, new systems often split up or abandon the kingdom, instead treating the protist groups as separate lines of eukaryotes. The recent scheme by Adl et al. (2005)[33] does not recognize formal ranks (phylum, class, etc.) and instead treats groups as clades of phylogenetically related organisms. This is intended to make the classification more stable in the long term and easier to update. Some of the main groups of protists, which may be treated as phyla, are listed in the taxobox, upper right.[42] Many are thought to be monophyletic, though there is still uncertainty. For instance, the Excavata are probably not monophyletic and the chromalveolates are probably only monophyletic if the haptophytes and cryptomonads are excluded.[43]
|
33 |
+
|
34 |
+
Nutrition can vary according to the type of protist. Most eukaryotic algae are autotrophic, but the pigments were lost in some groups.[vague] Other protists are heterotrophic, and may present phagotrophy, osmotrophy, saprotrophy or parasitism. Some are mixotrophic. Some protists that do not have / lost chloroplasts/mitochondria have entered into endosymbiontic relationship with other bacteria/algae to replace the missing functionality. For example, Paramecium bursaria and Paulinella have captured a green alga (Zoochlorella) and a cyanobacterium respectively that act as replacements for chloroplast. Meanwhile, a protist, Mixotricha paradoxa that has lost its mitochondria uses endosymbiontic bacteria as mitochondria and ectosymbiontic hair-like bacteria (Treponema spirochetes) for locomotion.
|
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+
|
36 |
+
Many protists are flagellate, for example, and filter feeding can take place where flagellates find prey. Other protists can engulf bacteria and other food particles, by extending their cell membrane around them to form a food vacuole and digesting them internally in a process termed phagocytosis.
|
37 |
+
|
38 |
+
For most important cellular structures and functions of animal and plants, it can be found a heritage among protists.[44]
|
39 |
+
|
40 |
+
Some protists reproduce sexually using gametes, while others reproduce asexually by binary fission.
|
41 |
+
|
42 |
+
Some species, for example Plasmodium falciparum, have extremely complex life cycles that involve multiple forms of the organism, some of which reproduce sexually and others asexually.[45] However, it is unclear how frequently sexual reproduction causes genetic exchange between different strains of Plasmodium in nature and most populations of parasitic protists may be clonal lines that rarely exchange genes with other members of their species.[46]
|
43 |
+
|
44 |
+
Eukaryotes emerged in evolution more than 1.5 billion years ago.[47] The earliest eukaryotes were likely protists. Although sexual reproduction is widespread among extant eukaryotes, it seemed unlikely until recently, that sex could be a primordial and fundamental characteristic of eukaryotes. A principal reason for this view was that sex appeared to be lacking in certain pathogenic protists whose ancestors branched off early from the eukaryotic family tree. However, several of these protists are now known to be capable of, or to recently have had the capability for, meiosis and hence sexual reproduction. For example, the common intestinal parasite Giardia lamblia was once considered to be a descendant of a protist lineage that predated the emergence of meiosis and sex. However, G. lamblia was recently found to have a core set of genes that function in meiosis and that are widely present among sexual eukaryotes.[48] These results suggested that G. lamblia is capable of meiosis and thus sexual reproduction. Furthermore, direct evidence for meiotic recombination, indicative of sex, was also found in G. lamblia.[49]
|
45 |
+
|
46 |
+
The pathogenic parasitic protists of the genus Leishmania have been shown to be capable of a sexual cycle in the invertebrate vector, likened to the meiosis undertaken in the trypanosomes.[50]
|
47 |
+
|
48 |
+
Trichomonas vaginalis, a parasitic protist, is not known to undergo meiosis, but when Malik et al.[51] tested for 29 genes that function in meiosis, they found 27 to be present, including 8 of 9 genes specific to meiosis in model eukaryotes. These findings suggest that T. vaginalis may be capable of meiosis. Since 21 of the 29 meiotic genes were also present in G. lamblia, it appears that most of these meiotic genes were likely present in a common ancestor of T. vaginalis and G. lamblia. These two species are descendants of protist lineages that are highly divergent among eukaryotes, leading Malik et al.[51] to suggest that these meiotic genes were likely present in a common ancestor of all eukaryotes.
|
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+
|
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+
Based on a phylogenetic analysis, Dacks and Roger proposed that facultative sex was present in the common ancestor of all eukaryotes.[52]
|
51 |
+
|
52 |
+
This view was further supported by a study of amoebae by Lahr et al.[53] Amoeba have generally been regarded as asexual protists. However, these authors describe evidence that most amoeboid lineages are anciently sexual, and that the majority of asexual groups likely arose recently and independently. Early researchers (e.g., Calkins) have interpreted phenomena related to chromidia (chromatin granules free in the cytoplasm) in amoeboid organisms as sexual reproduction.[54]
|
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+
|
54 |
+
Protists generally reproduce asexually under favorable environmental conditions, but tend to reproduce sexually under stressful conditions, such as starvation or heat shock.[55] Oxidative stress, which is associated with the production of reactive oxygen species leading to DNA damage, also appears to be an important factor in the induction of sex in protists.[55]
|
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+
|
56 |
+
Some commonly found Protist pathogens such as Toxoplasma gondii are capable of infecting and undergoing asexual reproduction in a wide variety of animals – which act as secondary or intermediate host – but can undergo sexual reproduction only in the primary or definitive host (for example: felids such as domestic cats in this case).[56][57][58]
|
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+
Free-living Protists occupy almost any environment that contains liquid water. Many protists, such as algae, are photosynthetic and are vital primary producers in ecosystems, particularly in the ocean as part of the plankton. Protists make up a large portion of the biomass in both marine and terrestrial environments.[59]
|
59 |
+
|
60 |
+
Other protists include pathogenic species, such as the kinetoplastid Trypanosoma brucei, which causes sleeping sickness, and species of the apicomplexan Plasmodium, which cause malaria.
|
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+
|
62 |
+
Some protists are significant parasites of animals (e.g.; five species of the parasitic genus Plasmodium cause malaria in humans and many others cause similar diseases in other vertebrates), plants[60][61] (the oomycete Phytophthora infestans causes late blight in potatoes)[62] or even of other protists.[63][64] Protist pathogens share many metabolic pathways with their eukaryotic hosts. This makes therapeutic target development extremely difficult – a drug that harms a protist parasite is also likely to harm its animal/plant host. A more thorough understanding of protist biology may allow these diseases to be treated more efficiently. For example, the apicoplast (a nonphotosynthetic chloroplast but essential to carry out important functions other than photosynthesis) present in apicomplexans provides an attractive target for treating diseases caused by dangerous pathogens such as plasmodium.
|
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+
Recent papers have proposed the use of viruses to treat infections caused by protozoa.[65][66]
|
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+
|
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+
Researchers from the Agricultural Research Service are taking advantage of protists as pathogens to control red imported fire ant (Solenopsis invicta) populations in Argentina. Spore-producing protists such as Kneallhazia solenopsae (recognized as a sister clade or the closest relative to the fungus kingdom now)[67] can reduce red fire ant populations by 53–100%.[68] Researchers have also been able to infect phorid fly parasitoids of the ant with the protist without harming the flies. This turns the flies into a vector that can spread the pathogenic protist between red fire ant colonies.[69]
|
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+
|
68 |
+
Many protists have neither hard parts nor resistant spores, and their fossils are extremely rare or unknown. Examples of such groups include the apicomplexans,[70] most ciliates,[71] some green algae (the Klebsormidiales),[72] choanoflagellates,[73] oomycetes,[74] brown algae,[75] yellow-green algae,[76] Excavata (e.g., euglenids).[77] Some of these have been found preserved in amber (fossilized tree resin) or under unusual conditions (e.g., Paleoleishmania, a kinetoplastid).
|
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+
|
70 |
+
Others are relatively common in the fossil record,[78] as the diatoms,[79] golden algae,[80] haptophytes (coccoliths),[81] silicoflagellates, tintinnids (ciliates), dinoflagellates,[82] green algae,[83] red algae,[84] heliozoans, radiolarians,[85] foraminiferans,[86] ebriids and testate amoebae (euglyphids, arcellaceans).[87] Some are even used as paleoecological indicators to reconstruct ancient environments.
|
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+
|
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+
More probable eukaryote fossils begin to appear at about 1.8 billion years ago, the acritarchs, spherical fossils of likely algal protists.[88] Another possible representative of early fossil eukaryotes are the Gabonionta.
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Supergroups[1] and typical phyla
|
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+
|
3 |
+
Many others;classification varies
|
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+
|
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+
A protist (/ˈproʊtɪst/) is any eukaryotic organism (one with cells containing a nucleus) that is not an animal, plant, or fungus. While it is likely that protists share a common ancestor (the last eukaryotic common ancestor),[2] the exclusion of other eukaryotes means that protists do not form a natural group, or clade.[a] So some protists may be more closely related to animals, plants, or fungi than they are to other protists; however, like algae, invertebrates, or protozoans, the grouping is used for convenience. The study of protists is termed protistology.[3]
|
6 |
+
|
7 |
+
The classification of a separate kingdom to the animals and plants was first proposed by John Hogg in 1860 as the kingdom Protoctista; in 1866 Ernst Haeckel also proposed a third kingdom Protista as "the kingdom of primitive forms".[4] Originally these also included prokaryotes, but with time these would be removed to a fourth kingdom Monera.[b] In the popular five-kingdom scheme proposed by Robert Whittaker in 1969, the Protista was defined as eukaryotic "organisms which are unicellular or unicellular-colonial and which form no tissues", and the fifth kingdom Fungi was established.[5][6][c] In the five-kingdom system of Lynn Margulis, the term protist is reserved for microscopic organisms, while the more inclusive kingdom Protoctista (or protoctists) included certain large multicellular eukaryotes, such as kelp, red algae and slime molds.[9] Others use the term protist interchangeably with Margulis's protoctist, to encompass both single-celled and multicellular eukaryotes, including those that form specialized tissues but do not fit into any of the other traditional kingdoms.[10]
|
8 |
+
|
9 |
+
Besides their relatively simple levels of organization, protists do not necessarily have much in common.[11] When used, the term "protists" is now considered to mean a paraphyletic assemblage of similar-appearing but diverse taxa (biological groups); these taxa do not have an exclusive common ancestor beyond being composed of eukaryotes, and have different life cycles, trophic levels, modes of locomotion and cellular structures.[12][13] Examples of protists include:[14] amoebas (including nucleariids and Foraminifera); choanaflagellates; ciliates; diatoms; dinoflagellates; Giardia; Plasmodium (which causes malaria); oomycetes (including Phytophthora, the cause of the Great Famine of Ireland); and slime molds. These examples are unicellular, although oomycetes can form filaments, and slime molds can aggregate.
|
10 |
+
|
11 |
+
In cladistic systems (classifications based on common ancestry), there are no equivalents to the taxa Protista or Protoctista, as both terms refer to a paraphyletic group that spans the entire eukaryotic tree of life. In cladistic classification, the contents of Protista are mostly distributed among various supergroups: examples include the SAR supergroup (of stramenopiles or heterokonts, alveolates, and Rhizaria); Archaeplastida (or Plantae sensu lato); Excavata (which is mostly unicellular flagellates); and Opisthokonta (which commonly includes unicellular flagellates, but also animals and fungi). "Protista", "Protoctista", and "Protozoa" are therefore considered obsolete. However, the term "protist" continues to be used informally as a catch-all term for eukayotic organisms that aren't within other traditional kingdoms. For example, the word "protist pathogen" may be used to denote any disease-causing organism that is not plant, animal, fungal, prokaryotic, viral, or subviral.[15]
|
12 |
+
|
13 |
+
The term protista was first used by Ernst Haeckel in 1866. Protists were traditionally subdivided into several groups based on similarities to the "higher" kingdoms such as:[4]
|
14 |
+
|
15 |
+
Some protists, sometimes called ambiregnal protists, have been considered to be both protozoa and algae or fungi (e.g., slime molds and flagellated algae), and names for these have been published under either or both of the ICN and the ICZN.[16][17] Conflicts, such as these – for example the dual-classification of Euglenids and Dinobryons, which are mixotrophic – is an example of why the kingdom Protista was adopted.
|
16 |
+
|
17 |
+
These traditional subdivisions, largely based on superficial commonalities, have been replaced by classifications based on phylogenetics (evolutionary relatedness among organisms).
|
18 |
+
Molecular analyses in modern taxonomy have been used to redistribute former members of this group into diverse and sometimes distantly related phyla. For instance, the water molds are now considered to be closely related to photosynthetic organisms such as Brown algae and Diatoms, the slime molds are grouped mainly under Amoebozoa, and the Amoebozoa itself includes only a subset of the "Amoeba" group, and significant number of erstwhile "Amoeboid" genera are distributed among Rhizarians and other Phyla.
|
19 |
+
|
20 |
+
However, the older terms are still used as informal names to describe the morphology and ecology of various protists. For example, the term protozoa is used to refer to heterotrophic species of protists that do not form filaments.
|
21 |
+
|
22 |
+
Among the pioneers in the study of the protists, which were almost ignored by Linnaeus except for some genera (e.g., Vorticella, Chaos, Volvox, Corallina, Conferva, Ulva, Chara, Fucus)[18][19] were Leeuwenhoek, O. F. Müller, C. G. Ehrenberg and Félix Dujardin.[20] The first groups used to classify microscopic organism were the Animalcules and the Infusoria.[21] In 1818, the German naturalist Georg August Goldfuss introduced the word Protozoa to refer to organisms such as ciliates and corals.[22][4] After the cell theory of Schwann and Schleiden (1838–39), this group was modified in 1848 by Carl von Siebold to include only animal-like unicellular organisms, such as foraminifera and amoebae.[23] The formal taxonomic category Protoctista was first proposed in the early 1860s by John Hogg, who argued that the protists should include what he saw as primitive unicellular forms of both plants and animals. He defined the Protoctista as a "fourth kingdom of nature", in addition to the then-traditional kingdoms of plants, animals and minerals.[24][4] The kingdom of minerals was later removed from taxonomy in 1866 by Ernst Haeckel, leaving plants, animals, and the protists (Protista), defined as a "kingdom of primitive forms".[25][26]
|
23 |
+
|
24 |
+
In 1938, Herbert Copeland resurrected Hogg's label, arguing that Haeckel's term Protista included anucleated microbes such as bacteria, which the term "Protoctista" (literally meaning "first established beings") did not. In contrast, Copeland's term included nucleated eukaryotes such as diatoms, green algae and fungi.[27] This classification was the basis for Whittaker's later definition of Fungi, Animalia, Plantae and Protista as the four kingdoms of life.[7] The kingdom Protista was later modified to separate prokaryotes into the separate kingdom of Monera, leaving the protists as a group of eukaryotic microorganisms.[5] These five kingdoms remained the accepted classification until the development of molecular phylogenetics in the late 20th century, when it became apparent that neither protists nor monera were single groups of related organisms (they were not monophyletic groups).[28]
|
25 |
+
|
26 |
+
Systematists today do not treat Protista as a formal taxon, but the term "protist" is still commonly used for convenience in two ways.[29] The most popular contemporary definition is a phylogenetic one, that identifies a paraphyletic group:[30] a protist is any eukaryote that is not an animal, (land) plant, or (true) fungus; this definition[31] excludes many unicellular groups, like the Microsporidia (fungi), many Chytridiomycetes (fungi), and yeasts (fungi), and also a non-unicellular group included in Protista in the past, the Myxozoa (animal).[32] Some systematists[who?] judge paraphyletic taxa acceptable, and use Protista in this sense as a formal taxon (as found in some secondary textbooks, for pedagogical purpose).[citation needed]
|
27 |
+
|
28 |
+
The other definition describes protists primarily by functional or biological criteria: protists are essentially those eukaryotes that are never multicellular,[29] that either exist as independent cells, or if they occur in colonies, do not show differentiation into tissues (but vegetative cell differentiation may occur restricted to sexual reproduction, alternate vegetative morphology, and quiescent or resistant stages, such as cysts);[33] this definition excludes many brown, multicellular red and green algae, which may have tissues.
|
29 |
+
|
30 |
+
The taxonomy of protists is still changing. Newer classifications attempt to present monophyletic groups based on morphological (especially ultrastructural),[34][35][36] biochemical (chemotaxonomy)[37][38] and DNA sequence (molecular research) information.[39][40] However, there are sometimes discordances between molecular and morphological investigations; these can be categorized as two types: (i) one morphology, multiple lineages (e.g. morphological convergence, cryptic species) and (ii) one lineage, multiple morphologies (e.g. phenotypic plasticity, multiple life-cycle stages).[41]
|
31 |
+
|
32 |
+
Because the protists as a whole are paraphyletic, new systems often split up or abandon the kingdom, instead treating the protist groups as separate lines of eukaryotes. The recent scheme by Adl et al. (2005)[33] does not recognize formal ranks (phylum, class, etc.) and instead treats groups as clades of phylogenetically related organisms. This is intended to make the classification more stable in the long term and easier to update. Some of the main groups of protists, which may be treated as phyla, are listed in the taxobox, upper right.[42] Many are thought to be monophyletic, though there is still uncertainty. For instance, the Excavata are probably not monophyletic and the chromalveolates are probably only monophyletic if the haptophytes and cryptomonads are excluded.[43]
|
33 |
+
|
34 |
+
Nutrition can vary according to the type of protist. Most eukaryotic algae are autotrophic, but the pigments were lost in some groups.[vague] Other protists are heterotrophic, and may present phagotrophy, osmotrophy, saprotrophy or parasitism. Some are mixotrophic. Some protists that do not have / lost chloroplasts/mitochondria have entered into endosymbiontic relationship with other bacteria/algae to replace the missing functionality. For example, Paramecium bursaria and Paulinella have captured a green alga (Zoochlorella) and a cyanobacterium respectively that act as replacements for chloroplast. Meanwhile, a protist, Mixotricha paradoxa that has lost its mitochondria uses endosymbiontic bacteria as mitochondria and ectosymbiontic hair-like bacteria (Treponema spirochetes) for locomotion.
|
35 |
+
|
36 |
+
Many protists are flagellate, for example, and filter feeding can take place where flagellates find prey. Other protists can engulf bacteria and other food particles, by extending their cell membrane around them to form a food vacuole and digesting them internally in a process termed phagocytosis.
|
37 |
+
|
38 |
+
For most important cellular structures and functions of animal and plants, it can be found a heritage among protists.[44]
|
39 |
+
|
40 |
+
Some protists reproduce sexually using gametes, while others reproduce asexually by binary fission.
|
41 |
+
|
42 |
+
Some species, for example Plasmodium falciparum, have extremely complex life cycles that involve multiple forms of the organism, some of which reproduce sexually and others asexually.[45] However, it is unclear how frequently sexual reproduction causes genetic exchange between different strains of Plasmodium in nature and most populations of parasitic protists may be clonal lines that rarely exchange genes with other members of their species.[46]
|
43 |
+
|
44 |
+
Eukaryotes emerged in evolution more than 1.5 billion years ago.[47] The earliest eukaryotes were likely protists. Although sexual reproduction is widespread among extant eukaryotes, it seemed unlikely until recently, that sex could be a primordial and fundamental characteristic of eukaryotes. A principal reason for this view was that sex appeared to be lacking in certain pathogenic protists whose ancestors branched off early from the eukaryotic family tree. However, several of these protists are now known to be capable of, or to recently have had the capability for, meiosis and hence sexual reproduction. For example, the common intestinal parasite Giardia lamblia was once considered to be a descendant of a protist lineage that predated the emergence of meiosis and sex. However, G. lamblia was recently found to have a core set of genes that function in meiosis and that are widely present among sexual eukaryotes.[48] These results suggested that G. lamblia is capable of meiosis and thus sexual reproduction. Furthermore, direct evidence for meiotic recombination, indicative of sex, was also found in G. lamblia.[49]
|
45 |
+
|
46 |
+
The pathogenic parasitic protists of the genus Leishmania have been shown to be capable of a sexual cycle in the invertebrate vector, likened to the meiosis undertaken in the trypanosomes.[50]
|
47 |
+
|
48 |
+
Trichomonas vaginalis, a parasitic protist, is not known to undergo meiosis, but when Malik et al.[51] tested for 29 genes that function in meiosis, they found 27 to be present, including 8 of 9 genes specific to meiosis in model eukaryotes. These findings suggest that T. vaginalis may be capable of meiosis. Since 21 of the 29 meiotic genes were also present in G. lamblia, it appears that most of these meiotic genes were likely present in a common ancestor of T. vaginalis and G. lamblia. These two species are descendants of protist lineages that are highly divergent among eukaryotes, leading Malik et al.[51] to suggest that these meiotic genes were likely present in a common ancestor of all eukaryotes.
|
49 |
+
|
50 |
+
Based on a phylogenetic analysis, Dacks and Roger proposed that facultative sex was present in the common ancestor of all eukaryotes.[52]
|
51 |
+
|
52 |
+
This view was further supported by a study of amoebae by Lahr et al.[53] Amoeba have generally been regarded as asexual protists. However, these authors describe evidence that most amoeboid lineages are anciently sexual, and that the majority of asexual groups likely arose recently and independently. Early researchers (e.g., Calkins) have interpreted phenomena related to chromidia (chromatin granules free in the cytoplasm) in amoeboid organisms as sexual reproduction.[54]
|
53 |
+
|
54 |
+
Protists generally reproduce asexually under favorable environmental conditions, but tend to reproduce sexually under stressful conditions, such as starvation or heat shock.[55] Oxidative stress, which is associated with the production of reactive oxygen species leading to DNA damage, also appears to be an important factor in the induction of sex in protists.[55]
|
55 |
+
|
56 |
+
Some commonly found Protist pathogens such as Toxoplasma gondii are capable of infecting and undergoing asexual reproduction in a wide variety of animals – which act as secondary or intermediate host – but can undergo sexual reproduction only in the primary or definitive host (for example: felids such as domestic cats in this case).[56][57][58]
|
57 |
+
|
58 |
+
Free-living Protists occupy almost any environment that contains liquid water. Many protists, such as algae, are photosynthetic and are vital primary producers in ecosystems, particularly in the ocean as part of the plankton. Protists make up a large portion of the biomass in both marine and terrestrial environments.[59]
|
59 |
+
|
60 |
+
Other protists include pathogenic species, such as the kinetoplastid Trypanosoma brucei, which causes sleeping sickness, and species of the apicomplexan Plasmodium, which cause malaria.
|
61 |
+
|
62 |
+
Some protists are significant parasites of animals (e.g.; five species of the parasitic genus Plasmodium cause malaria in humans and many others cause similar diseases in other vertebrates), plants[60][61] (the oomycete Phytophthora infestans causes late blight in potatoes)[62] or even of other protists.[63][64] Protist pathogens share many metabolic pathways with their eukaryotic hosts. This makes therapeutic target development extremely difficult – a drug that harms a protist parasite is also likely to harm its animal/plant host. A more thorough understanding of protist biology may allow these diseases to be treated more efficiently. For example, the apicoplast (a nonphotosynthetic chloroplast but essential to carry out important functions other than photosynthesis) present in apicomplexans provides an attractive target for treating diseases caused by dangerous pathogens such as plasmodium.
|
63 |
+
|
64 |
+
Recent papers have proposed the use of viruses to treat infections caused by protozoa.[65][66]
|
65 |
+
|
66 |
+
Researchers from the Agricultural Research Service are taking advantage of protists as pathogens to control red imported fire ant (Solenopsis invicta) populations in Argentina. Spore-producing protists such as Kneallhazia solenopsae (recognized as a sister clade or the closest relative to the fungus kingdom now)[67] can reduce red fire ant populations by 53–100%.[68] Researchers have also been able to infect phorid fly parasitoids of the ant with the protist without harming the flies. This turns the flies into a vector that can spread the pathogenic protist between red fire ant colonies.[69]
|
67 |
+
|
68 |
+
Many protists have neither hard parts nor resistant spores, and their fossils are extremely rare or unknown. Examples of such groups include the apicomplexans,[70] most ciliates,[71] some green algae (the Klebsormidiales),[72] choanoflagellates,[73] oomycetes,[74] brown algae,[75] yellow-green algae,[76] Excavata (e.g., euglenids).[77] Some of these have been found preserved in amber (fossilized tree resin) or under unusual conditions (e.g., Paleoleishmania, a kinetoplastid).
|
69 |
+
|
70 |
+
Others are relatively common in the fossil record,[78] as the diatoms,[79] golden algae,[80] haptophytes (coccoliths),[81] silicoflagellates, tintinnids (ciliates), dinoflagellates,[82] green algae,[83] red algae,[84] heliozoans, radiolarians,[85] foraminiferans,[86] ebriids and testate amoebae (euglyphids, arcellaceans).[87] Some are even used as paleoecological indicators to reconstruct ancient environments.
|
71 |
+
|
72 |
+
More probable eukaryote fossils begin to appear at about 1.8 billion years ago, the acritarchs, spherical fossils of likely algal protists.[88] Another possible representative of early fossil eukaryotes are the Gabonionta.
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1 |
+
Supergroups[1] and typical phyla
|
2 |
+
|
3 |
+
Many others;classification varies
|
4 |
+
|
5 |
+
A protist (/ˈproʊtɪst/) is any eukaryotic organism (one with cells containing a nucleus) that is not an animal, plant, or fungus. While it is likely that protists share a common ancestor (the last eukaryotic common ancestor),[2] the exclusion of other eukaryotes means that protists do not form a natural group, or clade.[a] So some protists may be more closely related to animals, plants, or fungi than they are to other protists; however, like algae, invertebrates, or protozoans, the grouping is used for convenience. The study of protists is termed protistology.[3]
|
6 |
+
|
7 |
+
The classification of a separate kingdom to the animals and plants was first proposed by John Hogg in 1860 as the kingdom Protoctista; in 1866 Ernst Haeckel also proposed a third kingdom Protista as "the kingdom of primitive forms".[4] Originally these also included prokaryotes, but with time these would be removed to a fourth kingdom Monera.[b] In the popular five-kingdom scheme proposed by Robert Whittaker in 1969, the Protista was defined as eukaryotic "organisms which are unicellular or unicellular-colonial and which form no tissues", and the fifth kingdom Fungi was established.[5][6][c] In the five-kingdom system of Lynn Margulis, the term protist is reserved for microscopic organisms, while the more inclusive kingdom Protoctista (or protoctists) included certain large multicellular eukaryotes, such as kelp, red algae and slime molds.[9] Others use the term protist interchangeably with Margulis's protoctist, to encompass both single-celled and multicellular eukaryotes, including those that form specialized tissues but do not fit into any of the other traditional kingdoms.[10]
|
8 |
+
|
9 |
+
Besides their relatively simple levels of organization, protists do not necessarily have much in common.[11] When used, the term "protists" is now considered to mean a paraphyletic assemblage of similar-appearing but diverse taxa (biological groups); these taxa do not have an exclusive common ancestor beyond being composed of eukaryotes, and have different life cycles, trophic levels, modes of locomotion and cellular structures.[12][13] Examples of protists include:[14] amoebas (including nucleariids and Foraminifera); choanaflagellates; ciliates; diatoms; dinoflagellates; Giardia; Plasmodium (which causes malaria); oomycetes (including Phytophthora, the cause of the Great Famine of Ireland); and slime molds. These examples are unicellular, although oomycetes can form filaments, and slime molds can aggregate.
|
10 |
+
|
11 |
+
In cladistic systems (classifications based on common ancestry), there are no equivalents to the taxa Protista or Protoctista, as both terms refer to a paraphyletic group that spans the entire eukaryotic tree of life. In cladistic classification, the contents of Protista are mostly distributed among various supergroups: examples include the SAR supergroup (of stramenopiles or heterokonts, alveolates, and Rhizaria); Archaeplastida (or Plantae sensu lato); Excavata (which is mostly unicellular flagellates); and Opisthokonta (which commonly includes unicellular flagellates, but also animals and fungi). "Protista", "Protoctista", and "Protozoa" are therefore considered obsolete. However, the term "protist" continues to be used informally as a catch-all term for eukayotic organisms that aren't within other traditional kingdoms. For example, the word "protist pathogen" may be used to denote any disease-causing organism that is not plant, animal, fungal, prokaryotic, viral, or subviral.[15]
|
12 |
+
|
13 |
+
The term protista was first used by Ernst Haeckel in 1866. Protists were traditionally subdivided into several groups based on similarities to the "higher" kingdoms such as:[4]
|
14 |
+
|
15 |
+
Some protists, sometimes called ambiregnal protists, have been considered to be both protozoa and algae or fungi (e.g., slime molds and flagellated algae), and names for these have been published under either or both of the ICN and the ICZN.[16][17] Conflicts, such as these – for example the dual-classification of Euglenids and Dinobryons, which are mixotrophic – is an example of why the kingdom Protista was adopted.
|
16 |
+
|
17 |
+
These traditional subdivisions, largely based on superficial commonalities, have been replaced by classifications based on phylogenetics (evolutionary relatedness among organisms).
|
18 |
+
Molecular analyses in modern taxonomy have been used to redistribute former members of this group into diverse and sometimes distantly related phyla. For instance, the water molds are now considered to be closely related to photosynthetic organisms such as Brown algae and Diatoms, the slime molds are grouped mainly under Amoebozoa, and the Amoebozoa itself includes only a subset of the "Amoeba" group, and significant number of erstwhile "Amoeboid" genera are distributed among Rhizarians and other Phyla.
|
19 |
+
|
20 |
+
However, the older terms are still used as informal names to describe the morphology and ecology of various protists. For example, the term protozoa is used to refer to heterotrophic species of protists that do not form filaments.
|
21 |
+
|
22 |
+
Among the pioneers in the study of the protists, which were almost ignored by Linnaeus except for some genera (e.g., Vorticella, Chaos, Volvox, Corallina, Conferva, Ulva, Chara, Fucus)[18][19] were Leeuwenhoek, O. F. Müller, C. G. Ehrenberg and Félix Dujardin.[20] The first groups used to classify microscopic organism were the Animalcules and the Infusoria.[21] In 1818, the German naturalist Georg August Goldfuss introduced the word Protozoa to refer to organisms such as ciliates and corals.[22][4] After the cell theory of Schwann and Schleiden (1838–39), this group was modified in 1848 by Carl von Siebold to include only animal-like unicellular organisms, such as foraminifera and amoebae.[23] The formal taxonomic category Protoctista was first proposed in the early 1860s by John Hogg, who argued that the protists should include what he saw as primitive unicellular forms of both plants and animals. He defined the Protoctista as a "fourth kingdom of nature", in addition to the then-traditional kingdoms of plants, animals and minerals.[24][4] The kingdom of minerals was later removed from taxonomy in 1866 by Ernst Haeckel, leaving plants, animals, and the protists (Protista), defined as a "kingdom of primitive forms".[25][26]
|
23 |
+
|
24 |
+
In 1938, Herbert Copeland resurrected Hogg's label, arguing that Haeckel's term Protista included anucleated microbes such as bacteria, which the term "Protoctista" (literally meaning "first established beings") did not. In contrast, Copeland's term included nucleated eukaryotes such as diatoms, green algae and fungi.[27] This classification was the basis for Whittaker's later definition of Fungi, Animalia, Plantae and Protista as the four kingdoms of life.[7] The kingdom Protista was later modified to separate prokaryotes into the separate kingdom of Monera, leaving the protists as a group of eukaryotic microorganisms.[5] These five kingdoms remained the accepted classification until the development of molecular phylogenetics in the late 20th century, when it became apparent that neither protists nor monera were single groups of related organisms (they were not monophyletic groups).[28]
|
25 |
+
|
26 |
+
Systematists today do not treat Protista as a formal taxon, but the term "protist" is still commonly used for convenience in two ways.[29] The most popular contemporary definition is a phylogenetic one, that identifies a paraphyletic group:[30] a protist is any eukaryote that is not an animal, (land) plant, or (true) fungus; this definition[31] excludes many unicellular groups, like the Microsporidia (fungi), many Chytridiomycetes (fungi), and yeasts (fungi), and also a non-unicellular group included in Protista in the past, the Myxozoa (animal).[32] Some systematists[who?] judge paraphyletic taxa acceptable, and use Protista in this sense as a formal taxon (as found in some secondary textbooks, for pedagogical purpose).[citation needed]
|
27 |
+
|
28 |
+
The other definition describes protists primarily by functional or biological criteria: protists are essentially those eukaryotes that are never multicellular,[29] that either exist as independent cells, or if they occur in colonies, do not show differentiation into tissues (but vegetative cell differentiation may occur restricted to sexual reproduction, alternate vegetative morphology, and quiescent or resistant stages, such as cysts);[33] this definition excludes many brown, multicellular red and green algae, which may have tissues.
|
29 |
+
|
30 |
+
The taxonomy of protists is still changing. Newer classifications attempt to present monophyletic groups based on morphological (especially ultrastructural),[34][35][36] biochemical (chemotaxonomy)[37][38] and DNA sequence (molecular research) information.[39][40] However, there are sometimes discordances between molecular and morphological investigations; these can be categorized as two types: (i) one morphology, multiple lineages (e.g. morphological convergence, cryptic species) and (ii) one lineage, multiple morphologies (e.g. phenotypic plasticity, multiple life-cycle stages).[41]
|
31 |
+
|
32 |
+
Because the protists as a whole are paraphyletic, new systems often split up or abandon the kingdom, instead treating the protist groups as separate lines of eukaryotes. The recent scheme by Adl et al. (2005)[33] does not recognize formal ranks (phylum, class, etc.) and instead treats groups as clades of phylogenetically related organisms. This is intended to make the classification more stable in the long term and easier to update. Some of the main groups of protists, which may be treated as phyla, are listed in the taxobox, upper right.[42] Many are thought to be monophyletic, though there is still uncertainty. For instance, the Excavata are probably not monophyletic and the chromalveolates are probably only monophyletic if the haptophytes and cryptomonads are excluded.[43]
|
33 |
+
|
34 |
+
Nutrition can vary according to the type of protist. Most eukaryotic algae are autotrophic, but the pigments were lost in some groups.[vague] Other protists are heterotrophic, and may present phagotrophy, osmotrophy, saprotrophy or parasitism. Some are mixotrophic. Some protists that do not have / lost chloroplasts/mitochondria have entered into endosymbiontic relationship with other bacteria/algae to replace the missing functionality. For example, Paramecium bursaria and Paulinella have captured a green alga (Zoochlorella) and a cyanobacterium respectively that act as replacements for chloroplast. Meanwhile, a protist, Mixotricha paradoxa that has lost its mitochondria uses endosymbiontic bacteria as mitochondria and ectosymbiontic hair-like bacteria (Treponema spirochetes) for locomotion.
|
35 |
+
|
36 |
+
Many protists are flagellate, for example, and filter feeding can take place where flagellates find prey. Other protists can engulf bacteria and other food particles, by extending their cell membrane around them to form a food vacuole and digesting them internally in a process termed phagocytosis.
|
37 |
+
|
38 |
+
For most important cellular structures and functions of animal and plants, it can be found a heritage among protists.[44]
|
39 |
+
|
40 |
+
Some protists reproduce sexually using gametes, while others reproduce asexually by binary fission.
|
41 |
+
|
42 |
+
Some species, for example Plasmodium falciparum, have extremely complex life cycles that involve multiple forms of the organism, some of which reproduce sexually and others asexually.[45] However, it is unclear how frequently sexual reproduction causes genetic exchange between different strains of Plasmodium in nature and most populations of parasitic protists may be clonal lines that rarely exchange genes with other members of their species.[46]
|
43 |
+
|
44 |
+
Eukaryotes emerged in evolution more than 1.5 billion years ago.[47] The earliest eukaryotes were likely protists. Although sexual reproduction is widespread among extant eukaryotes, it seemed unlikely until recently, that sex could be a primordial and fundamental characteristic of eukaryotes. A principal reason for this view was that sex appeared to be lacking in certain pathogenic protists whose ancestors branched off early from the eukaryotic family tree. However, several of these protists are now known to be capable of, or to recently have had the capability for, meiosis and hence sexual reproduction. For example, the common intestinal parasite Giardia lamblia was once considered to be a descendant of a protist lineage that predated the emergence of meiosis and sex. However, G. lamblia was recently found to have a core set of genes that function in meiosis and that are widely present among sexual eukaryotes.[48] These results suggested that G. lamblia is capable of meiosis and thus sexual reproduction. Furthermore, direct evidence for meiotic recombination, indicative of sex, was also found in G. lamblia.[49]
|
45 |
+
|
46 |
+
The pathogenic parasitic protists of the genus Leishmania have been shown to be capable of a sexual cycle in the invertebrate vector, likened to the meiosis undertaken in the trypanosomes.[50]
|
47 |
+
|
48 |
+
Trichomonas vaginalis, a parasitic protist, is not known to undergo meiosis, but when Malik et al.[51] tested for 29 genes that function in meiosis, they found 27 to be present, including 8 of 9 genes specific to meiosis in model eukaryotes. These findings suggest that T. vaginalis may be capable of meiosis. Since 21 of the 29 meiotic genes were also present in G. lamblia, it appears that most of these meiotic genes were likely present in a common ancestor of T. vaginalis and G. lamblia. These two species are descendants of protist lineages that are highly divergent among eukaryotes, leading Malik et al.[51] to suggest that these meiotic genes were likely present in a common ancestor of all eukaryotes.
|
49 |
+
|
50 |
+
Based on a phylogenetic analysis, Dacks and Roger proposed that facultative sex was present in the common ancestor of all eukaryotes.[52]
|
51 |
+
|
52 |
+
This view was further supported by a study of amoebae by Lahr et al.[53] Amoeba have generally been regarded as asexual protists. However, these authors describe evidence that most amoeboid lineages are anciently sexual, and that the majority of asexual groups likely arose recently and independently. Early researchers (e.g., Calkins) have interpreted phenomena related to chromidia (chromatin granules free in the cytoplasm) in amoeboid organisms as sexual reproduction.[54]
|
53 |
+
|
54 |
+
Protists generally reproduce asexually under favorable environmental conditions, but tend to reproduce sexually under stressful conditions, such as starvation or heat shock.[55] Oxidative stress, which is associated with the production of reactive oxygen species leading to DNA damage, also appears to be an important factor in the induction of sex in protists.[55]
|
55 |
+
|
56 |
+
Some commonly found Protist pathogens such as Toxoplasma gondii are capable of infecting and undergoing asexual reproduction in a wide variety of animals – which act as secondary or intermediate host – but can undergo sexual reproduction only in the primary or definitive host (for example: felids such as domestic cats in this case).[56][57][58]
|
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Free-living Protists occupy almost any environment that contains liquid water. Many protists, such as algae, are photosynthetic and are vital primary producers in ecosystems, particularly in the ocean as part of the plankton. Protists make up a large portion of the biomass in both marine and terrestrial environments.[59]
|
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+
|
60 |
+
Other protists include pathogenic species, such as the kinetoplastid Trypanosoma brucei, which causes sleeping sickness, and species of the apicomplexan Plasmodium, which cause malaria.
|
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+
|
62 |
+
Some protists are significant parasites of animals (e.g.; five species of the parasitic genus Plasmodium cause malaria in humans and many others cause similar diseases in other vertebrates), plants[60][61] (the oomycete Phytophthora infestans causes late blight in potatoes)[62] or even of other protists.[63][64] Protist pathogens share many metabolic pathways with their eukaryotic hosts. This makes therapeutic target development extremely difficult – a drug that harms a protist parasite is also likely to harm its animal/plant host. A more thorough understanding of protist biology may allow these diseases to be treated more efficiently. For example, the apicoplast (a nonphotosynthetic chloroplast but essential to carry out important functions other than photosynthesis) present in apicomplexans provides an attractive target for treating diseases caused by dangerous pathogens such as plasmodium.
|
63 |
+
|
64 |
+
Recent papers have proposed the use of viruses to treat infections caused by protozoa.[65][66]
|
65 |
+
|
66 |
+
Researchers from the Agricultural Research Service are taking advantage of protists as pathogens to control red imported fire ant (Solenopsis invicta) populations in Argentina. Spore-producing protists such as Kneallhazia solenopsae (recognized as a sister clade or the closest relative to the fungus kingdom now)[67] can reduce red fire ant populations by 53–100%.[68] Researchers have also been able to infect phorid fly parasitoids of the ant with the protist without harming the flies. This turns the flies into a vector that can spread the pathogenic protist between red fire ant colonies.[69]
|
67 |
+
|
68 |
+
Many protists have neither hard parts nor resistant spores, and their fossils are extremely rare or unknown. Examples of such groups include the apicomplexans,[70] most ciliates,[71] some green algae (the Klebsormidiales),[72] choanoflagellates,[73] oomycetes,[74] brown algae,[75] yellow-green algae,[76] Excavata (e.g., euglenids).[77] Some of these have been found preserved in amber (fossilized tree resin) or under unusual conditions (e.g., Paleoleishmania, a kinetoplastid).
|
69 |
+
|
70 |
+
Others are relatively common in the fossil record,[78] as the diatoms,[79] golden algae,[80] haptophytes (coccoliths),[81] silicoflagellates, tintinnids (ciliates), dinoflagellates,[82] green algae,[83] red algae,[84] heliozoans, radiolarians,[85] foraminiferans,[86] ebriids and testate amoebae (euglyphids, arcellaceans).[87] Some are even used as paleoecological indicators to reconstruct ancient environments.
|
71 |
+
|
72 |
+
More probable eukaryote fossils begin to appear at about 1.8 billion years ago, the acritarchs, spherical fossils of likely algal protists.[88] Another possible representative of early fossil eukaryotes are the Gabonionta.
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en/4839.html.txt
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Supergroups[1] and typical phyla
|
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+
|
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+
Many others;classification varies
|
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|
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+
A protist (/ˈproʊtɪst/) is any eukaryotic organism (one with cells containing a nucleus) that is not an animal, plant, or fungus. While it is likely that protists share a common ancestor (the last eukaryotic common ancestor),[2] the exclusion of other eukaryotes means that protists do not form a natural group, or clade.[a] So some protists may be more closely related to animals, plants, or fungi than they are to other protists; however, like algae, invertebrates, or protozoans, the grouping is used for convenience. The study of protists is termed protistology.[3]
|
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+
|
7 |
+
The classification of a separate kingdom to the animals and plants was first proposed by John Hogg in 1860 as the kingdom Protoctista; in 1866 Ernst Haeckel also proposed a third kingdom Protista as "the kingdom of primitive forms".[4] Originally these also included prokaryotes, but with time these would be removed to a fourth kingdom Monera.[b] In the popular five-kingdom scheme proposed by Robert Whittaker in 1969, the Protista was defined as eukaryotic "organisms which are unicellular or unicellular-colonial and which form no tissues", and the fifth kingdom Fungi was established.[5][6][c] In the five-kingdom system of Lynn Margulis, the term protist is reserved for microscopic organisms, while the more inclusive kingdom Protoctista (or protoctists) included certain large multicellular eukaryotes, such as kelp, red algae and slime molds.[9] Others use the term protist interchangeably with Margulis's protoctist, to encompass both single-celled and multicellular eukaryotes, including those that form specialized tissues but do not fit into any of the other traditional kingdoms.[10]
|
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+
|
9 |
+
Besides their relatively simple levels of organization, protists do not necessarily have much in common.[11] When used, the term "protists" is now considered to mean a paraphyletic assemblage of similar-appearing but diverse taxa (biological groups); these taxa do not have an exclusive common ancestor beyond being composed of eukaryotes, and have different life cycles, trophic levels, modes of locomotion and cellular structures.[12][13] Examples of protists include:[14] amoebas (including nucleariids and Foraminifera); choanaflagellates; ciliates; diatoms; dinoflagellates; Giardia; Plasmodium (which causes malaria); oomycetes (including Phytophthora, the cause of the Great Famine of Ireland); and slime molds. These examples are unicellular, although oomycetes can form filaments, and slime molds can aggregate.
|
10 |
+
|
11 |
+
In cladistic systems (classifications based on common ancestry), there are no equivalents to the taxa Protista or Protoctista, as both terms refer to a paraphyletic group that spans the entire eukaryotic tree of life. In cladistic classification, the contents of Protista are mostly distributed among various supergroups: examples include the SAR supergroup (of stramenopiles or heterokonts, alveolates, and Rhizaria); Archaeplastida (or Plantae sensu lato); Excavata (which is mostly unicellular flagellates); and Opisthokonta (which commonly includes unicellular flagellates, but also animals and fungi). "Protista", "Protoctista", and "Protozoa" are therefore considered obsolete. However, the term "protist" continues to be used informally as a catch-all term for eukayotic organisms that aren't within other traditional kingdoms. For example, the word "protist pathogen" may be used to denote any disease-causing organism that is not plant, animal, fungal, prokaryotic, viral, or subviral.[15]
|
12 |
+
|
13 |
+
The term protista was first used by Ernst Haeckel in 1866. Protists were traditionally subdivided into several groups based on similarities to the "higher" kingdoms such as:[4]
|
14 |
+
|
15 |
+
Some protists, sometimes called ambiregnal protists, have been considered to be both protozoa and algae or fungi (e.g., slime molds and flagellated algae), and names for these have been published under either or both of the ICN and the ICZN.[16][17] Conflicts, such as these – for example the dual-classification of Euglenids and Dinobryons, which are mixotrophic – is an example of why the kingdom Protista was adopted.
|
16 |
+
|
17 |
+
These traditional subdivisions, largely based on superficial commonalities, have been replaced by classifications based on phylogenetics (evolutionary relatedness among organisms).
|
18 |
+
Molecular analyses in modern taxonomy have been used to redistribute former members of this group into diverse and sometimes distantly related phyla. For instance, the water molds are now considered to be closely related to photosynthetic organisms such as Brown algae and Diatoms, the slime molds are grouped mainly under Amoebozoa, and the Amoebozoa itself includes only a subset of the "Amoeba" group, and significant number of erstwhile "Amoeboid" genera are distributed among Rhizarians and other Phyla.
|
19 |
+
|
20 |
+
However, the older terms are still used as informal names to describe the morphology and ecology of various protists. For example, the term protozoa is used to refer to heterotrophic species of protists that do not form filaments.
|
21 |
+
|
22 |
+
Among the pioneers in the study of the protists, which were almost ignored by Linnaeus except for some genera (e.g., Vorticella, Chaos, Volvox, Corallina, Conferva, Ulva, Chara, Fucus)[18][19] were Leeuwenhoek, O. F. Müller, C. G. Ehrenberg and Félix Dujardin.[20] The first groups used to classify microscopic organism were the Animalcules and the Infusoria.[21] In 1818, the German naturalist Georg August Goldfuss introduced the word Protozoa to refer to organisms such as ciliates and corals.[22][4] After the cell theory of Schwann and Schleiden (1838–39), this group was modified in 1848 by Carl von Siebold to include only animal-like unicellular organisms, such as foraminifera and amoebae.[23] The formal taxonomic category Protoctista was first proposed in the early 1860s by John Hogg, who argued that the protists should include what he saw as primitive unicellular forms of both plants and animals. He defined the Protoctista as a "fourth kingdom of nature", in addition to the then-traditional kingdoms of plants, animals and minerals.[24][4] The kingdom of minerals was later removed from taxonomy in 1866 by Ernst Haeckel, leaving plants, animals, and the protists (Protista), defined as a "kingdom of primitive forms".[25][26]
|
23 |
+
|
24 |
+
In 1938, Herbert Copeland resurrected Hogg's label, arguing that Haeckel's term Protista included anucleated microbes such as bacteria, which the term "Protoctista" (literally meaning "first established beings") did not. In contrast, Copeland's term included nucleated eukaryotes such as diatoms, green algae and fungi.[27] This classification was the basis for Whittaker's later definition of Fungi, Animalia, Plantae and Protista as the four kingdoms of life.[7] The kingdom Protista was later modified to separate prokaryotes into the separate kingdom of Monera, leaving the protists as a group of eukaryotic microorganisms.[5] These five kingdoms remained the accepted classification until the development of molecular phylogenetics in the late 20th century, when it became apparent that neither protists nor monera were single groups of related organisms (they were not monophyletic groups).[28]
|
25 |
+
|
26 |
+
Systematists today do not treat Protista as a formal taxon, but the term "protist" is still commonly used for convenience in two ways.[29] The most popular contemporary definition is a phylogenetic one, that identifies a paraphyletic group:[30] a protist is any eukaryote that is not an animal, (land) plant, or (true) fungus; this definition[31] excludes many unicellular groups, like the Microsporidia (fungi), many Chytridiomycetes (fungi), and yeasts (fungi), and also a non-unicellular group included in Protista in the past, the Myxozoa (animal).[32] Some systematists[who?] judge paraphyletic taxa acceptable, and use Protista in this sense as a formal taxon (as found in some secondary textbooks, for pedagogical purpose).[citation needed]
|
27 |
+
|
28 |
+
The other definition describes protists primarily by functional or biological criteria: protists are essentially those eukaryotes that are never multicellular,[29] that either exist as independent cells, or if they occur in colonies, do not show differentiation into tissues (but vegetative cell differentiation may occur restricted to sexual reproduction, alternate vegetative morphology, and quiescent or resistant stages, such as cysts);[33] this definition excludes many brown, multicellular red and green algae, which may have tissues.
|
29 |
+
|
30 |
+
The taxonomy of protists is still changing. Newer classifications attempt to present monophyletic groups based on morphological (especially ultrastructural),[34][35][36] biochemical (chemotaxonomy)[37][38] and DNA sequence (molecular research) information.[39][40] However, there are sometimes discordances between molecular and morphological investigations; these can be categorized as two types: (i) one morphology, multiple lineages (e.g. morphological convergence, cryptic species) and (ii) one lineage, multiple morphologies (e.g. phenotypic plasticity, multiple life-cycle stages).[41]
|
31 |
+
|
32 |
+
Because the protists as a whole are paraphyletic, new systems often split up or abandon the kingdom, instead treating the protist groups as separate lines of eukaryotes. The recent scheme by Adl et al. (2005)[33] does not recognize formal ranks (phylum, class, etc.) and instead treats groups as clades of phylogenetically related organisms. This is intended to make the classification more stable in the long term and easier to update. Some of the main groups of protists, which may be treated as phyla, are listed in the taxobox, upper right.[42] Many are thought to be monophyletic, though there is still uncertainty. For instance, the Excavata are probably not monophyletic and the chromalveolates are probably only monophyletic if the haptophytes and cryptomonads are excluded.[43]
|
33 |
+
|
34 |
+
Nutrition can vary according to the type of protist. Most eukaryotic algae are autotrophic, but the pigments were lost in some groups.[vague] Other protists are heterotrophic, and may present phagotrophy, osmotrophy, saprotrophy or parasitism. Some are mixotrophic. Some protists that do not have / lost chloroplasts/mitochondria have entered into endosymbiontic relationship with other bacteria/algae to replace the missing functionality. For example, Paramecium bursaria and Paulinella have captured a green alga (Zoochlorella) and a cyanobacterium respectively that act as replacements for chloroplast. Meanwhile, a protist, Mixotricha paradoxa that has lost its mitochondria uses endosymbiontic bacteria as mitochondria and ectosymbiontic hair-like bacteria (Treponema spirochetes) for locomotion.
|
35 |
+
|
36 |
+
Many protists are flagellate, for example, and filter feeding can take place where flagellates find prey. Other protists can engulf bacteria and other food particles, by extending their cell membrane around them to form a food vacuole and digesting them internally in a process termed phagocytosis.
|
37 |
+
|
38 |
+
For most important cellular structures and functions of animal and plants, it can be found a heritage among protists.[44]
|
39 |
+
|
40 |
+
Some protists reproduce sexually using gametes, while others reproduce asexually by binary fission.
|
41 |
+
|
42 |
+
Some species, for example Plasmodium falciparum, have extremely complex life cycles that involve multiple forms of the organism, some of which reproduce sexually and others asexually.[45] However, it is unclear how frequently sexual reproduction causes genetic exchange between different strains of Plasmodium in nature and most populations of parasitic protists may be clonal lines that rarely exchange genes with other members of their species.[46]
|
43 |
+
|
44 |
+
Eukaryotes emerged in evolution more than 1.5 billion years ago.[47] The earliest eukaryotes were likely protists. Although sexual reproduction is widespread among extant eukaryotes, it seemed unlikely until recently, that sex could be a primordial and fundamental characteristic of eukaryotes. A principal reason for this view was that sex appeared to be lacking in certain pathogenic protists whose ancestors branched off early from the eukaryotic family tree. However, several of these protists are now known to be capable of, or to recently have had the capability for, meiosis and hence sexual reproduction. For example, the common intestinal parasite Giardia lamblia was once considered to be a descendant of a protist lineage that predated the emergence of meiosis and sex. However, G. lamblia was recently found to have a core set of genes that function in meiosis and that are widely present among sexual eukaryotes.[48] These results suggested that G. lamblia is capable of meiosis and thus sexual reproduction. Furthermore, direct evidence for meiotic recombination, indicative of sex, was also found in G. lamblia.[49]
|
45 |
+
|
46 |
+
The pathogenic parasitic protists of the genus Leishmania have been shown to be capable of a sexual cycle in the invertebrate vector, likened to the meiosis undertaken in the trypanosomes.[50]
|
47 |
+
|
48 |
+
Trichomonas vaginalis, a parasitic protist, is not known to undergo meiosis, but when Malik et al.[51] tested for 29 genes that function in meiosis, they found 27 to be present, including 8 of 9 genes specific to meiosis in model eukaryotes. These findings suggest that T. vaginalis may be capable of meiosis. Since 21 of the 29 meiotic genes were also present in G. lamblia, it appears that most of these meiotic genes were likely present in a common ancestor of T. vaginalis and G. lamblia. These two species are descendants of protist lineages that are highly divergent among eukaryotes, leading Malik et al.[51] to suggest that these meiotic genes were likely present in a common ancestor of all eukaryotes.
|
49 |
+
|
50 |
+
Based on a phylogenetic analysis, Dacks and Roger proposed that facultative sex was present in the common ancestor of all eukaryotes.[52]
|
51 |
+
|
52 |
+
This view was further supported by a study of amoebae by Lahr et al.[53] Amoeba have generally been regarded as asexual protists. However, these authors describe evidence that most amoeboid lineages are anciently sexual, and that the majority of asexual groups likely arose recently and independently. Early researchers (e.g., Calkins) have interpreted phenomena related to chromidia (chromatin granules free in the cytoplasm) in amoeboid organisms as sexual reproduction.[54]
|
53 |
+
|
54 |
+
Protists generally reproduce asexually under favorable environmental conditions, but tend to reproduce sexually under stressful conditions, such as starvation or heat shock.[55] Oxidative stress, which is associated with the production of reactive oxygen species leading to DNA damage, also appears to be an important factor in the induction of sex in protists.[55]
|
55 |
+
|
56 |
+
Some commonly found Protist pathogens such as Toxoplasma gondii are capable of infecting and undergoing asexual reproduction in a wide variety of animals – which act as secondary or intermediate host – but can undergo sexual reproduction only in the primary or definitive host (for example: felids such as domestic cats in this case).[56][57][58]
|
57 |
+
|
58 |
+
Free-living Protists occupy almost any environment that contains liquid water. Many protists, such as algae, are photosynthetic and are vital primary producers in ecosystems, particularly in the ocean as part of the plankton. Protists make up a large portion of the biomass in both marine and terrestrial environments.[59]
|
59 |
+
|
60 |
+
Other protists include pathogenic species, such as the kinetoplastid Trypanosoma brucei, which causes sleeping sickness, and species of the apicomplexan Plasmodium, which cause malaria.
|
61 |
+
|
62 |
+
Some protists are significant parasites of animals (e.g.; five species of the parasitic genus Plasmodium cause malaria in humans and many others cause similar diseases in other vertebrates), plants[60][61] (the oomycete Phytophthora infestans causes late blight in potatoes)[62] or even of other protists.[63][64] Protist pathogens share many metabolic pathways with their eukaryotic hosts. This makes therapeutic target development extremely difficult – a drug that harms a protist parasite is also likely to harm its animal/plant host. A more thorough understanding of protist biology may allow these diseases to be treated more efficiently. For example, the apicoplast (a nonphotosynthetic chloroplast but essential to carry out important functions other than photosynthesis) present in apicomplexans provides an attractive target for treating diseases caused by dangerous pathogens such as plasmodium.
|
63 |
+
|
64 |
+
Recent papers have proposed the use of viruses to treat infections caused by protozoa.[65][66]
|
65 |
+
|
66 |
+
Researchers from the Agricultural Research Service are taking advantage of protists as pathogens to control red imported fire ant (Solenopsis invicta) populations in Argentina. Spore-producing protists such as Kneallhazia solenopsae (recognized as a sister clade or the closest relative to the fungus kingdom now)[67] can reduce red fire ant populations by 53–100%.[68] Researchers have also been able to infect phorid fly parasitoids of the ant with the protist without harming the flies. This turns the flies into a vector that can spread the pathogenic protist between red fire ant colonies.[69]
|
67 |
+
|
68 |
+
Many protists have neither hard parts nor resistant spores, and their fossils are extremely rare or unknown. Examples of such groups include the apicomplexans,[70] most ciliates,[71] some green algae (the Klebsormidiales),[72] choanoflagellates,[73] oomycetes,[74] brown algae,[75] yellow-green algae,[76] Excavata (e.g., euglenids).[77] Some of these have been found preserved in amber (fossilized tree resin) or under unusual conditions (e.g., Paleoleishmania, a kinetoplastid).
|
69 |
+
|
70 |
+
Others are relatively common in the fossil record,[78] as the diatoms,[79] golden algae,[80] haptophytes (coccoliths),[81] silicoflagellates, tintinnids (ciliates), dinoflagellates,[82] green algae,[83] red algae,[84] heliozoans, radiolarians,[85] foraminiferans,[86] ebriids and testate amoebae (euglyphids, arcellaceans).[87] Some are even used as paleoecological indicators to reconstruct ancient environments.
|
71 |
+
|
72 |
+
More probable eukaryote fossils begin to appear at about 1.8 billion years ago, the acritarchs, spherical fossils of likely algal protists.[88] Another possible representative of early fossil eukaryotes are the Gabonionta.
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An avalanche (also called a snowslide) is an event that occurs when a cohesive slab of snow lying upon a weaker layer of snow fractures and slides down a steep slope. Avalanches are typically triggered in a starting zone from a mechanical failure in the snowpack (slab avalanche) when the forces of the snow exceed its strength but sometimes only with gradual widening (loose snow avalanche). After initiation, avalanches usually accelerate rapidly and grow in mass and volume as they entrain more snow. If the avalanche moves fast enough, some of the snow may mix with the air forming a powder snow avalanche, which is a type of gravity current.
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|
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Slides of rocks or debris, behaving in a similar way to snow, are also referred to as avalanches (see rockslide[1]). The remainder of this article refers to snow avalanches, because they are the ones that occur the most in the world.
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|
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The load on the snowpack may be only due to gravity, in which case failure may result either from weakening in the snowpack or increased load due to precipitation. Avalanches initiated by this process are known as spontaneous avalanches. Avalanches can also be triggered by other loading conditions such as human or biologically related activities. Seismic activity (activities in the earth's crust relating to plate movements) may also trigger the failure in the snowpack and avalanches.
|
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|
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Although primarily composed of flowing snow and air, large avalanches have the capability to entrain ice, rocks, trees, and other surficial material. However, they are distinct from slushflows which have higher water content and more laminar flow, mudslides which have greater fluidity, rock slides which are often ice free, and serac collapses during an icefall. Avalanches are not rare or random events and are endemic to any mountain range that accumulates a standing snowpack. Avalanches are most common during winter or spring but glacier movements may cause ice and snow avalanches at any time of year. In mountainous terrain, avalanches are among the most serious objective natural hazards to life and property, with their destructive capability resulting from their potential to carry enormous masses of snow at high speeds.
|
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|
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+
There is no universally accepted classification system for different forms of avalanches. Avalanches can be described by their size, their destructive potential, their initiation mechanism, their composition and their dynamics.
|
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|
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+
Most avalanches occur spontaneously during storms under increased load due to snowfall and/or erosion. The second largest cause of natural avalanches is metamorphic changes in the snowpack such as melting due to solar radiation. Other natural causes include rain, earthquakes, rockfall and icefall. Artificial triggers of avalanches include skiers, snowmobiles, and controlled explosive work. Contrary to popular belief, avalanches are not triggered by loud sound; the pressure from sound is orders of magnitude too small to trigger an avalanche.[2]
|
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|
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+
Avalanche initiation can start at a point with only a small amount of snow moving initially; this is typical of wet snow avalanches or avalanches in dry unconsolidated snow. However, if the snow has sintered into a stiff slab overlying a weak layer then fractures can propagate very rapidly, so that a large volume of snow, that may be thousands of cubic meters, can start moving almost simultaneously.
|
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|
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+
A snowpack will fail when the load exceeds the strength. The load is straightforward; it is the weight of the snow. However, the strength of the snowpack is much more difficult to determine and is extremely heterogeneous. It varies in detail with properties of the snow grains, size, density, morphology, temperature, water content; and the properties of the bonds between the grains.
|
16 |
+
[3] These properties may all metamorphose in time according to the local humidity, water vapour flux, temperature and heat flux. The top of the snowpack is also extensively influenced by incoming radiation and the local air flow. One of the aims of avalanche research is to develop and validate computer models that can describe the evolution of the seasonal snowpack over time.[4] A complicating factor is the complex interaction of terrain and weather, which causes significant spatial and temporal variability of the depths, crystal forms, and layering of the seasonal snowpack.
|
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+
|
18 |
+
Slab avalanches form frequently in snow that has been deposited, or redeposited by wind. They have the characteristic appearance of a block (slab) of snow cut out from its surroundings by fractures. Elements of slab avalanches include the following: a crown fracture at the top of the start zone, flank fractures on the sides of the start zones, and a fracture at the bottom called the stauchwall. The crown and flank fractures are vertical walls in the snow delineating the snow that was entrained in the avalanche from the snow that remained on the slope. Slabs can vary in thickness from a few centimetres to three metres. Slab avalanches account for around 90% of avalanche-related fatalities in backcountry users.
|
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+
|
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+
The largest avalanches form turbulent suspension currents known as powder snow avalanches or mixed avalanches.[5] These consist of a powder cloud, which overlies a dense avalanche. They can form from any type of snow or initiation mechanism, but usually occur with fresh dry powder.
|
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+
They can exceed speeds of 300 kilometres per hour (190 mph), and masses of 10000000 tonnes; their flows can travel long distances along flat valley bottoms and even uphill for short distances.
|
22 |
+
|
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+
In contrast to powder snow avalanches, wet snow avalanches are a low velocity suspension of snow and water, with the flow confined to the track surface (McClung, first edition 1999, page 108).[3] The low speed of travel is due to the friction between the sliding surface of the track and the water saturated flow. Despite the low speed of travel (~10–40 km/h), wet snow avalanches are capable of generating powerful destructive forces, due to the large mass and density. The body of the flow of a wet snow avalanche can plough through soft snow, and can scour boulders, earth, trees, and other vegetation; leaving exposed and often scored ground in the avalanche track. Wet snow avalanches can be initiated from either loose snow releases, or slab releases, and only occur in snow packs that are water saturated and isothermally equilibrated to the melting point of water. The isothermal characteristic of wet snow avalanches has led to the secondary term of isothermal slides found in the literature (for example in Daffern, 1999, page 93).[6] At temperate latitudes wet snow avalanches are frequently associated with climatic avalanche cycles at the end of the winter season, when there is significant daytime warming.
|
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+
|
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+
As an avalanche moves down a slope it follows a certain pathway that is dependent on the slope's degree of steepness and the volume of snow/ice involved in the mass movement. The origin of an avalanche is called the Starting Point and typically occurs on a 30–45 degree slope. The body of the pathway is called the Track of the avalanche and usually occurs on a 20–30 degree slope. When the avalanche loses its momentum and eventually stops it reaches the Runout Zone. This usually occurs when the slope has reached a steepness that is less than 20 degrees.[7] These degrees are not consistently true due to the fact that each avalanche is unique depending on the stability of the snowpack that it was derived from as well as the environmental or human influences that triggered the mass movement.
|
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+
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+
People caught in avalanches can die from suffocation, trauma, or hypothermia. On average, 28 people die in avalanches every winter in the United States.[8] Globally, an average of over 150 people die each year from avalanches.
|
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+
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+
An ice avalanche occurs when a large piece of ice, such as from a serac or calving glacier, falls onto ice (such as the Khumbu Icefall), triggering a movement of broken ice chunks. The resulting movement is more analogous to a rockfall or a landslide than a snow avalanche.[3] They are typically very difficult to predict and almost impossible to mitigate.
|
30 |
+
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+
Doug Fesler and Jill Fredston developed a conceptual model of the three primary elements of avalanches: terrain, weather, and snowpack. Terrain describes the places where avalanches occur, weather describes the meteorological conditions that create the snowpack, and snowpack describes the structural characteristics of snow that make avalanche formation possible.[3][9]
|
32 |
+
|
33 |
+
Avalanche formation requires a slope shallow enough for snow to accumulate but steep enough for the snow to accelerate once set in motion by the combination of mechanical failure (of the snowpack) and gravity. The angle of the slope that can hold snow, called the angle of repose, depends on a variety of factors such as crystal form and moisture content. Some forms of drier and colder snow will only stick to shallower slopes, while wet and warm snow can bond to very steep surfaces. In particular, in coastal mountains, such as the Cordillera del Paine region of Patagonia, deep snowpacks collect on vertical and even overhanging rock faces. The slope angle that can allow moving snow to accelerate depends on a variety of factors such as the snow's shear strength (which is itself dependent upon crystal form) and the configuration of layers and inter-layer interfaces.
|
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+
|
35 |
+
The snowpack on slopes with sunny exposures is strongly influenced by sunshine. Diurnal cycles of thawing and refreezing can stabilize the snowpack by promoting settlement. Strong freeze-thaw cycles result in the formation of surface crusts during the night and of unstable surface snow during the day. Slopes in the lee of a ridge or of another wind obstacle accumulate more snow and are more likely to include pockets of deep snow, wind slabs, and cornices, all of which, when disturbed, may result in avalanche formation. Conversely, the snowpack on a windward slope is often much shallower than on a lee slope.
|
36 |
+
|
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+
Avalanches and avalanche paths share common elements: a start zone where the avalanche originates, a track along which the avalanche flows, and a runout zone where the avalanche comes to rest. The debris deposit is the accumulated mass of the avalanched snow once it has come to rest in the runout zone. For the image at left, many small avalanches form in this avalanche path every year, but most of these avalanches do not run the full vertical or horizontal length of the path. The frequency with which avalanches form in a given area is known as the return period.
|
38 |
+
|
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+
The start zone of an avalanche must be steep enough to allow snow to accelerate once set in motion, additionally convex slopes are less stable than concave slopes, because of the disparity between the tensile strength of snow layers and their compressive strength. The composition and structure of the ground surface beneath the snowpack influences the stability of the snowpack, either being a source of strength or weakness. Avalanches are unlikely to form in very thick forests, but boulders and sparsely distributed vegetation can create weak areas deep within the snowpack through the formation of strong temperature gradients. Full-depth avalanches (avalanches that sweep a slope virtually clean of snow cover) are more common on slopes with smooth ground, such as grass or rock slabs.
|
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+
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+
Generally speaking, avalanches follow drainages down-slope, frequently sharing drainage features with summertime watersheds. At and below tree line, avalanche paths through drainages are well defined by vegetation boundaries called trim lines, which occur where avalanches have removed trees and prevented regrowth of large vegetation. Engineered drainages, such as the avalanche dam on Mount Stephen in Kicking Horse Pass, have been constructed to protect people and property by redirecting the flow of avalanches. Deep debris deposits from avalanches will collect in catchments at the terminus of a run out, such as gullies and river beds.
|
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+
|
43 |
+
Slopes flatter than 25 degrees or steeper than 60 degrees typically have a lower incidence of avalanches. Human-triggered avalanches have the greatest incidence when the snow's angle of repose is between 35 and 45 degrees; the critical angle, the angle at which human-triggered avalanches are most frequent, is 38 degrees. When the incidence of human triggered avalanches is normalized by the rates of recreational use, however, hazard increases uniformly with slope angle, and no significant difference in hazard for a given exposure direction can be found.[10] The rule of thumb is: A slope that is flat enough to hold snow but steep enough to ski has the potential to generate an avalanche, regardless of the angle.
|
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+
|
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The snowpack is composed of ground-parallel layers that accumulate over the winter. Each layer contains ice grains that are representative of the distinct meteorological conditions during which the snow formed and was deposited. Once deposited, a snow layer continues to evolve under the influence of the meteorological conditions that prevail after deposition.
|
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+
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For an avalanche to occur, it is necessary that a snowpack have a weak layer (or instability) below a slab of cohesive snow. In practice the formal mechanical and structural factors related to snowpack instability are not directly observable outside of laboratories, thus the more easily observed properties of the snow layers (e.g. penetration resistance, grain size, grain type, temperature) are used as index measurements of the mechanical properties of the snow (e.g. tensile strength, friction coefficients, shear strength, and ductile strength). This results in two principal sources of uncertainty in determining snowpack stability based on snow structure: First, both the factors influencing snow stability and the specific characteristics of the snowpack vary widely within small areas and time scales, resulting in significant difficulty extrapolating point observations of snow layers across different scales of space and time. Second, the relationship between readily observable snowpack characteristics and the snowpack's critical mechanical properties has not been completely developed.
|
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|
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+
While the deterministic relationship between snowpack characteristics and snowpack stability is still a matter of ongoing scientific study, there is a growing empirical understanding of the snow composition and deposition characteristics that influence the likelihood of an avalanche. Observation and experience has shown that newly fallen snow requires time to bond with the snow layers beneath it, especially if the new snow falls during very cold and dry conditions. If ambient air temperatures are cold enough, shallow snow above or around boulders, plants, and other discontinuities in the slope, weakens from rapid crystal growth that occurs in the presence of a critical temperature gradient. Large, angular snow crystals are indicators of weak snow, because such crystals have fewer bonds per unit volume than small, rounded crystals that pack tightly together. Consolidated snow is less likely to slough than loose powdery layers or wet isothermal snow; however, consolidated snow is a necessary condition for the occurrence of slab avalanches, and persistent instabilities within the snowpack can hide below well-consolidated surface layers. Uncertainty associated with the empirical understanding of the factors influencing snow stability leads most professional avalanche workers to recommend conservative use of avalanche terrain relative to current snowpack instability.
|
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+
|
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Avalanches can only occur in a standing snowpack. Typically winter seasons at high latitudes, high altitudes, or both have weather that is sufficiently unsettled and cold enough for precipitated snow to accumulate into a seasonal snowpack. Continentality, through its potentiating influence on the meteorological extremes experienced by snowpacks, is an important factor in the evolution of instabilities, and consequential occurrence of avalanches. Conversely, proximity to coastal environments moderates the meteorological extremes experienced by snowpacks, and results in a faster stabilization of the snowpack after storm cycles.[11] The evolution of the snowpack is critically sensitive to small variations within the narrow range of meteorological conditions that allow for the accumulation of snow into a snowpack. Among the critical factors controlling snowpack evolution are: heating by the sun, radiational cooling, vertical temperature gradients in standing snow, snowfall amounts, and snow types. Generally, mild winter weather will promote the settlement and stabilization of the snowpack; conversely, very cold, windy, or hot weather will weaken the snowpack.
|
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+
|
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+
At temperatures close to the freezing point of water, or during times of moderate solar radiation, a gentle freeze-thaw cycle will take place. The melting and refreezing of water in the snow strengthens the snowpack during the freezing phase and weakens it during the thawing phase. A rapid rise in temperature, to a point significantly above the freezing point of water, may cause avalanche formation at any time of year.
|
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+
|
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+
Persistent cold temperatures can either prevent new snow from stabilizing or destabilize the existing snowpack. Cold air temperatures on the snow surface produce a temperature gradient in the snow, because the ground temperature at the base of the snowpack is usually around 0 °C, and the ambient air temperature can be much colder. When a temperature gradient greater than 10 °C change per vertical meter of snow is sustained for more than a day, angular crystals called depth hoar or facets begin forming in the snowpack because of rapid moisture transport along the temperature gradient. These angular crystals, which bond poorly to one another and the surrounding snow, often become a persistent weakness in the snowpack. When a slab lying on top of a persistent weakness is loaded by a force greater than the strength of the slab and persistent weak layer, the persistent weak layer can fail and generate an avalanche.
|
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Any wind stronger than a light breeze can contribute to a rapid accumulation of snow on sheltered slopes downwind. Wind slab forms quickly and, if present, weaker snow below the slab may not have time to adjust to the new load. Even on a clear day, wind can quickly load a slope with snow by blowing snow from one place to another. Top-loading occurs when wind deposits snow from the top of a slope; cross-loading occurs when wind deposits snow parallel to the slope. When a wind blows over the top of a mountain, the leeward, or downwind, side of the mountain experiences top-loading, from the top to the bottom of that lee slope. When the wind blows across a ridge that leads up the mountain, the leeward side of the ridge is subject to cross-loading. Cross-loaded wind-slabs are usually difficult to identify visually.
|
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+
|
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+
Snowstorms and rainstorms are important contributors to avalanche danger. Heavy snowfall will cause instability in the existing snowpack, both because of the additional weight and because the new snow has insufficient time to bond to underlying snow layers. Rain has a similar effect. In the short-term, rain causes instability because, like a heavy snowfall, it imposes an additional load on the snowpack; and, once rainwater seeps down through the snow, it acts as a lubricant, reducing the natural friction between snow layers that holds the snowpack together. Most avalanches happen during or soon after a storm.
|
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+
|
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+
Daytime exposure to sunlight will rapidly destabilize the upper layers of the snowpack if the sunlight is strong enough to melt the snow, thereby reducing its hardness. During clear nights, the snowpack can re-freeze when ambient air temperatures fall below freezing, through the process of long-wave radiative cooling, or both. Radiative heat loss occurs when the night air is significantly cooler than the snowpack, and the heat stored in the snow is re-radiated into the atmosphere.
|
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|
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+
When a slab avalanche forms, the slab disintegrates into increasingly smaller fragments as the snow travels downhill. If the fragments become small enough the outer layer of the avalanche, called a saltation layer, takes on the characteristics of a fluid. When sufficiently fine particles are present they can become airborne and, given a sufficient quantity of airborne snow, this portion of the avalanche can become separated from the bulk of the avalanche and travel a greater distance as a powder snow avalanche.[12] Scientific studies using radar, following the 1999 Galtür avalanche disaster, confirmed the hypothesis that a saltation layer forms between the surface and the airborne components of an avalanche, which can also separate from the bulk of the avalanche.[13]
|
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+
|
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+
Driving an avalanche is the component of the avalanche's weight parallel to the slope; as the avalanche progresses any unstable snow in its path will tend to become incorporated, so increasing the overall weight. This force will increase as the steepness of the slope increases, and diminish as the slope flattens. Resisting this are a number of components that are thought to interact with each other: the friction between the avalanche and the surface beneath; friction between the air and snow within the fluid; fluid-dynamic drag at the leading edge of the avalanche; shear resistance between the avalanche and the air through which it is passing, and shear resistance between the fragments within the avalanche itself. An avalanche will continue to accelerate until the resistance exceeds the forward force.[14]
|
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|
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+
Attempts to model avalanche behaviour date from the early 20th century, notably the work of Professor Lagotala in preparation for the 1924 Winter Olympics in Chamonix.[15] His method was developed by A. Voellmy and popularised following the publication in 1955 of his Ueber die Zerstoerungskraft von Lawinen (On the Destructive Force of Avalanches).[16]
|
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+
|
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+
Voellmy used a simple empirical formula, treating an avalanche as a sliding block of snow moving with a drag force that was proportional to the square of the speed of its flow:[17]
|
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+
|
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+
He and others subsequently derived other formulae that take other factors into account, with the Voellmy-Salm-Gubler and the Perla-Cheng-McClung models becoming most widely used as simple tools to model flowing (as opposed to powder snow) avalanches.[15]
|
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+
|
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Since the 1990s many more sophisticated models have been developed. In Europe much of the recent work was carried out as part of the SATSIE (Avalanche Studies and Model Validation in Europe) research project supported by the European Commission[18] which produced the leading-edge MN2L model, now in use with the Service Restauration des Terrains en Montagne (Mountain Rescue Service) in
|
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+
France, and D2FRAM (Dynamical Two-Flow-Regime Avalanche Model), which was still undergoing validation as of 2007.[19] Other known models are the SAMOS-AT avalanche simulation software[20] and the RAMMS software.[21]
|
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|
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Preventative measures are employed in areas where avalanches pose a significant threat to people, such as ski resorts, mountain towns, roads, and railways. There are several ways to prevent avalanches and lessen their power and develop preventative measures to reduce the likelihood and size of avalanches by disrupting the structure of the snowpack, while passive measures reinforce and stabilize the snowpack in situ. The simplest active measure is repeatedly traveling on a snowpack as snow accumulates; this can be by means of boot-packing, ski-cutting, or machine grooming. Explosives are used extensively to prevent avalanches, by triggering smaller avalanches that break down instabilities in the snowpack, and removing overburden that can result in larger avalanches. Explosive charges are delivered by a number of methods including hand-tossed charges, helicopter-dropped bombs, Gazex concussion lines, and ballistic projectiles launched by air cannons and artillery. Passive preventive systems such as snow fences and light walls can be used to direct the placement of snow. Snow builds up around the fence, especially the side that faces the prevailing winds. Downwind of the fence, snow buildup is lessened. This is caused by the loss of snow at the fence that would have been deposited and the pickup of the snow that is already there by the wind, which was depleted of snow at the fence. When there is a sufficient density of trees, they can greatly reduce the strength of avalanches. They hold snow in place and when there is an avalanche, the impact of the snow against the trees slows it down. Trees can either be planted or they can be conserved, such as in the building of a ski resort, to reduce the strength of avalanches.
|
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|
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In turn, socio-environmental changes can influence the occurrence of damaging avalanches: some studies linking changes in land-use/land-cover patterns and the evolution of snow avalanche damage in mid latitude mountains show the importance of the role played by vegetation cover, that is at the root of the increase of damage when the protective forest is deforested (because of demographic growth, intensive grazing and industrial or legal causes), and at the root of the decrease of damage because of the transformation of a traditional land-management system based on overexploitation into a system based on land marginalization and reforestation, something that has happened mainly since the mid-20th century in mountain environments of developed countries[22]
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|
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In many areas, regular avalanche tracks can be identified and precautions can be taken to minimise damage, such as the prevention of development in these areas. To mitigate the effect of avalanches the construction of artificial barriers can be very effective in reducing avalanche damage. There are several types: One kind of barrier (snow net) uses a net strung between poles that are anchored by guy wires in addition to their foundations. These barriers are similar to those used for rockslides. Another type of barrier is a rigid fence-like structure (snow fence) and may be constructed of steel, wood or pre-stressed concrete. They usually have gaps between the beams and are built perpendicular to the slope, with reinforcing beams on the downhill side. Rigid barriers are often considered unsightly, especially when many rows must be built. They are also expensive and vulnerable to damage from falling rocks in the warmer months. In addition to industrially manufactured barriers, landscaped barriers, called avalanche dams stop or deflect avalanches with their weight and strength. These barriers are made out of concrete, rocks or earth. They are usually placed right above the structure, road or railway that they are trying to protect, although they can also be used to channel avalanches into other barriers. Occasionally, earth mounds are placed in the avalanche's path to slow it down. Finally, along transportation corridors, large shelters, called snow sheds, can be built directly in the slide path of an avalanche to protect traffic from avalanches.
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Warning systems can detect avalanches which develop slowly, such as ice avalanches caused by icefalls from glaciers. Interferometric Radars, high-resolution Cameras, or motion sensors can monitor instable areas over a long term, lasting from days to years. Experts interpret the recorded data and are able to recognize upcoming ruptures in order to initiate appropriate measures. Such systems (e.g. the monitoring of the Weissmies glacier in Switzerland[23]) can recognize events several days in advance.
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Modern radar technology enables the monitoring of large areas and the localization of avalanches at any weather condition, by day and by night. Complex alarm systems are able to detect avalanches within a short time in order to close (e.g. roads and rails) or evacuate (e.g. construction sites) endangered areas. An example of such a system is installed on the only access road of Zermatt in Switzerland.[24] Two radars monitor the slope of a mountain above the road. The system automatically closes the road by activating several barriers and traffic lights within seconds such that no people are harmed.
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+
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Avalanche accidents are broadly differentiated into 2 categories: accidents in recreational settings, and accidents in residential, industrial, and transportation settings. This distinction is motivated by the observed difference in the causes of avalanche accidents in the two settings. In the recreational setting most accidents are caused by the people involved in the avalanche. In a 1996 study, Jamieson et al. (pages 7–20)[25] found that 83% of all avalanches in the recreational setting were caused by those who were involved in the accident. In contrast, all of the accidents in the residential, industrial, and transportation settings were due to spontaneous natural avalanches. Because of the difference in the causes of avalanche accidents, and the activities pursued in the two settings, avalanche and disaster management professionals have developed two related preparedness, rescue, and recovery strategies for each of the settings.
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Two avalanches occurred in March 1910 in the Cascade and Selkirk Mountain ranges; On March 1 the Wellington avalanche killed 96 in Washington state, United States. Three days later 62 railroad workers were killed in the Rogers Pass avalanche in British Columbia, Canada.
|
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+
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+
During World War I, an estimated 40,000 to 80,000 soldiers died as a result of avalanches during the mountain campaign in the Alps at the Austrian-Italian front, many of which were caused by artillery fire.[26][27] Some 10,000 men, from both sides, lost their lives in avalanches in December 1916.[28]
|
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+
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In the northern hemisphere winter of 1950–1951 approximately 649 avalanches were recorded in a three-month period throughout the Alps in Austria, France, Switzerland, Italy and Germany. This series of avalanches killed around 265 people and was termed the Winter of Terror.
|
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+
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A mountain climbing camp on Lenin Peak, in what is now Kyrgyzstan, was wiped out in 1990 when an earthquake triggered a large avalanche that overran the camp.[29] Forty-three climbers were killed.[30]
|
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+
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+
In 1993, the Bayburt Üzengili avalanche killed 60 individuals in Üzengili in the province of Bayburt, Turkey.
|
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+
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A large avalanche in Montroc, France, in 1999, 300,000 cubic metres of snow slid on a 30° slope, achieving a speed in the region of 100 km/h (62 mph). It killed 12 people in their chalets under 100,000 tons of snow, 5 meters (16 feet) deep. The mayor of Chamonix was convicted of second-degree murder for not evacuating the area, but received a suspended sentence.[31]
|
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+
|
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The small Austrian village of Galtür was hit by the Galtür avalanche in 1999. The village was thought to be in a safe zone but the avalanche was exceptionally large and flowed into the village. Thirty-one people died.
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+
|
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+
On December 1, 2000, the Glory Bowl Avalanche formed on Mt. Glory which is located within the Teton Mountain Range in Wyoming, United States. Joel Roof was snowboarding recreationally in this backcountry, bowl-shaped run and triggered the avalanche. He was carried nearly 2,000 feet to the base of the mountain and was not successfully rescued.[32]
|
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+
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In Europe, the avalanche risk is widely rated on the following scale, which was adopted in April 1993 to replace the earlier non-standard national schemes. Descriptions were last updated in May 2003 to enhance uniformity.[33]
|
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|
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+
In France, most avalanche deaths occur at risk levels 3 and 4. In Switzerland most occur at levels 2 and 3. It is thought that this may be due to national differences of interpretation when assessing the risks.[34]
|
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|
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[1] Stability:
|
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|
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[2] additional load:
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Gradient:
|
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Avalanche size:[citation needed]
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In the United States and Canada, the following avalanche danger scale is used. Descriptors vary depending on country.
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There are nine different types of avalanche problems:[35][36]
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The Canadian classification for avalanche size is based upon the consequences of the avalanche. Half sizes are commonly used.[37]
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The size of avalanches are classified using two scales; size relative to destructive force or D-scale and size relative to the avalanche path or R-scale.[38][39] Both size scales range from 1 to 5 with the D size scale half sizes can be used.[38][39]
|
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Slab avalanche hazard analysis can be done using the Rutschblock Test. A 2 m wide block of snow is isolated from the rest of the slope and progressively loaded. The result is a rating of slope stability on a seven step scale.[40]
|
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(Rutsch means slide in German).
|
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Climate change-caused temperature increases and changes in precipitation patterns will likely differ between the different mountain regions.[41] But overall a rising seasonal snow line and a decrease in the number of days with snow cover are predicted.[41][42] The impacts of these changes on avalanches are different at different elevations. At lower elevations a long-term reduction in the number of avalanches corresponding to a decrease in snow, and a short-term increase in the number of wet avalanches are predicted.[41][43][44][45] Higher elevations predicted to remain about the seasonal snow line, will likely see an increase in avalanche activity in regions seeing increases in precipitation during the winter season.[44][46] Predictions also show an increase in the number of rain on snow events,[42] and wet avalanche cycles occurring earlier in the spring during the remainder of this century.[47]
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Media related to Avalanche chute at Wikimedia Commons
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A proverb (from Latin: proverbium) is a simple, concrete, traditional saying that expresses a perceived truth based on common sense or experience. Proverbs are often metaphorical and use formulaic language. Collectively, they form a genre of folklore.
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Some proverbs exist in more than one language because people borrow them from languages and cultures similar to theirs. In the West, the Bible (including, but not limited to the Book of Proverbs) and medieval Latin (aided by the work of Erasmus) have played a considerable role in distributing proverbs. Not all Biblical proverbs, however, were distributed to the same extent: one scholar has gathered evidence to show that cultures in which the Bible is the "major spiritual book contain between three hundred and five hundred proverbs that stem from the Bible,"[1] whereas another shows that, of the 106 most common and widespread proverbs across Europe, eleven are from the Bible.[2] However, almost every culture has its own unique proverbs.
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Lord John Russell (c. 1850) observed poetically that a "proverb is the wit of one, and the wisdom of many."[3] But giving the word "proverb" the sort of definition theorists need has proven to be a difficult task, and although scholars often quote Archer Taylor's argument that formulating a scientific "definition of a proverb is too difficult to repay the undertaking... An incommunicable quality tells us this sentence is proverbial and that one is not. Hence no definition will enable us to identify positively a sentence as proverbial,"[4] many students of proverbs have attempted to itemize its essential characteristics.
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More constructively, Mieder has proposed the following definition, "A proverb is a short, generally known sentence of the folk which contains wisdom, truth, morals, and traditional views in a metaphorical, fixed, and memorizable form and which is handed down from generation to generation".[5] Norrick created a table of distinctive features to distinguish proverbs from idioms, cliches, etc.[6] Prahlad distinguishes proverbs from some other, closely related types of sayings, "True proverbs must further be distinguished from other types of proverbial speech, e.g. proverbial phrases, Wellerisms, maxims, quotations, and proverbial comparisons."[7] Based on Persian proverbs, Zolfaghari and Ameri propose the following definition: "A proverb is a short sentence, which is well-known and at times rhythmic, including advice, sage themes and ethnic experiences, comprising simile, metaphor or irony which is well-known among people for its fluent wording, clarity of expression, simplicity, expansiveness and generality and is used either with or without change."[8]
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There are many sayings in English that are commonly referred to as "proverbs", such as weather sayings. Alan Dundes, however, rejects including such sayings among truly proverbs: "Are weather proverbs proverbs? I would say emphatically 'No!'"[9] The definition of "proverb" has also changed over the years. For example, the following was labeled "A Yorkshire proverb" in 1883, but would not be categorized as a proverb by most today, "as throng as Throp's wife when she hanged herself with a dish-cloth".[10] The changing of the definition of "proverb" is also noted in Turkish.[11]
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In other languages and cultures, the definition of "proverb" also differs from English. In the Chumburung language of Ghana, "aŋase are literal proverbs and akpare are metaphoric ones".[12] Among the Bini of Nigeria, there are three words that are used to translate "proverb": ere, ivbe, and itan. The first relates to historical events, the second relates to current events, and the third was "linguistic ornamentation in formal discourse".[13] Among the Balochi of Pakistan and Afghanistan, there is a word batal for ordinary proverbs and bassīttuks for "proverbs with background stories".[14]
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There are also language communities that combine proverbs and riddles in some sayings, leading some scholars to create the label "proverb riddles".[15][16][17]
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Proverbs come from a variety of sources.[19] Some are, indeed, the result of people pondering and crafting language, such as some by Confucius, Plato, Baltasar Gracián, etc. Others are taken from such diverse sources as poetry,[20] stories,[21] songs, commercials, advertisements, movies, literature, etc.[22] A number of the well known sayings of Jesus, Shakespeare, and others have become proverbs, though they were original at the time of their creation, and many of these sayings were not seen as proverbs when they were first coined. Many proverbs are also based on stories, often the end of a story. For example, the proverb "Who will bell the cat?" is from the end of a story about the mice planning how to be safe from the cat.[23]
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Some authors have created proverbs in their writings, such as J.R.R. Tolkien,[24][25] and some of these proverbs have made their way into broader society. Similarly, C.S. Lewis' created proverb about a lobster in a pot, from the Chronicles of Narnia, has also gained currency.[26] In cases like this, deliberately created proverbs for fictional societies have become proverbs in real societies. In a fictional story set in a real society, the movie Forrest Gump introduced "Life is like a box of chocolates" into broad society.[27] In at least one case, it appears that a proverb deliberately created by one writer has been naively picked up and used by another who assumed it to be an established Chinese proverb, Ford Madox Ford having picked up a proverb from Ernest Bramah, "It would be hypocrisy to seek for the person of the Sacred Emperor in a Low Tea House."[28]
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The proverb with "a longer history than any other recorded proverb in the world", going back to "around 1800 BC"[29] is in a Sumerian clay tablet, "The bitch by her acting too hastily brought forth the blind".[30][31] Though many proverbs are ancient, they were all newly created at some point by somebody. Sometimes it is easy to detect that a proverb is newly coined by a reference to something recent, such as the Haitian proverb "The fish that is being microwaved doesn't fear the lightning".[32] Similarly, there is a recent Maltese proverb, wil-muturi, ferh u duluri "Women and motorcycles are joys and griefs"; the proverb is clearly new, but still formed as a traditional style couplet with rhyme.[33] Also, there is a proverb in the Kafa language of Ethiopia that refers to the forced military conscription of the 1980s, "...the one who hid himself lived to have children."[34] A Mongolian proverb also shows evidence of recent origin, "A beggar who sits on gold; Foam rubber piled on edge."[35] A political candidate in Kenya popularised a new proverb in his 1995 campaign, Chuth ber "Immediacy is best". "The proverb has since been used in other contexts to prompt quick action."[36] Over 1,400 new English proverbs are said to have been coined and gained currency in the 20th century.[37] This process of creating proverbs is always ongoing, so that possible new proverbs are being created constantly. Those sayings that are adopted and used by an adequate number of people become proverbs in that society.[38][39]
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Interpreting proverbs is often complex, but is best done in a context.[40] Interpreting proverbs from other cultures is much more difficult than interpreting proverbs in one's own culture. Even within English-speaking cultures, there is difference of opinion on how to interpret the proverb "A rolling stone gathers no moss." Some see it as condemning a person that keeps moving, seeing moss as a positive thing, such as profit; others see the proverb as praising people that keep moving and developing, seeing moss as a negative thing, such as negative habits.[41]
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Similarly, among Tajik speakers, the proverb "One hand cannot clap" has two significantly different interpretations. Most see the proverb as promoting teamwork. Others understand it to mean that an argument requires two people.[42] In an extreme example, one researcher working in Ghana found that for a single Akan proverb, twelve different interpretations were given.[43] Proveb interpretation is not automatic, even for people within a culture: Owomoyela tells of a Yoruba radio program that asked people to interpret an unfamiliar Yoruba proverb, "very few people could do so".[44] Siran found that people who had moved out of the traditional Vute-speaking area of Cameroon were not able to interpret Vute proverbs correctly, even though they still spoke Vute. Their interpretations tended to be literal.[45]
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Children will sometimes interpret proverbs in a literal sense, not yet knowing how to understand the conventionalized metaphor. Interpretation of proverbs is also affected by injuries and diseases of the brain, "A hallmark of schizophrenia is impaired proverb interpretation."[46]
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Proverbs in various languages are found with a wide variety of grammatical structures.[47] In English, for example, we find the following structures (in addition to others):
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However, people will often quote only a fraction of a proverb to invoke an entire proverb, e.g. "All is fair" instead of "All is fair in love and war", and "A rolling stone" for "A rolling stone gathers no moss."
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The grammar of proverbs is not always the typical grammar of the spoken language, often elements are moved around, to achieve rhyme or focus.[48]
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Another type of grammatical construction is the wellerism, a speaker and a quotation, often with an unusual circumstance, such as the following, a representative of a wellerism proverb found in many languages: "The bride couldn't dance; she said, 'The room floor isn't flat.'"[49]
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Another type of grammatical structure in proverbs is a short dialogue:
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Because many proverbs are both poetic and traditional, they are often passed down in fixed forms. Though spoken language may change, many proverbs are often preserved in conservative, even archaic, form. In English, for example, "betwixt" is not used by many, but a form of it is still heard (or read) in the proverb "There is many a slip 'twixt the cup and the lip." The conservative form preserves the meter and the rhyme. This conservative nature of proverbs can result in archaic words and grammatical structures being preserved in individual proverbs, as has been documented in Amharic,[54] Greek,[55] Nsenga,[56] Polish,[57] Venda[58] and Hebrew.[59]
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In addition, proverbs may still be used in languages which were once more widely known in a society, but are now no longer so widely known. For example, English speakers use some non-English proverbs that are drawn from languages that used to be widely understood by the educated class, e.g. "C'est la vie" from French and "Carpe diem" from Latin.
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Proverbs are often handed down through generations. Therefore, "many proverbs refer to old measurements, obscure professions, outdated weapons, unknown plants, animals, names, and various other traditional matters."[60]
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Therefore, it is common that they preserve words that become less common and archaic in broader society.[61][62] Proverbs in solid form—such as murals, carvings, and glass—can be viewed even after the language of their form is no longer widely understood, such as an Anglo-French proverb in a stained glass window in York.[63]
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Proverbs are often and easily translated and transferred from one language into another. "There is nothing so uncertain as the derivation of proverbs, the same proverb being often found in all nations, and it is impossible to assign its paternity."[64]
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Proverbs are often borrowed across lines of language, religion, and even time. For example, a proverb of the approximate form "No flies enter a mouth that is shut" is currently found in Spain, France, Ethiopia, and many countries in between. It is embraced as a true local proverb in many places and should not be excluded in any collection of proverbs because it is shared by the neighbors. However, though it has gone through multiple languages and millennia, the proverb can be traced back to an ancient Babylonian proverb (Pritchard 1958:146). Another example of a widely spread proverb is "A drowning person clutches at [frogs] foam", found in Peshai of Afghanistan[65] and Orma of Kenya,[66] and presumably places in between.
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Proverbs about one hand clapping are common across Asia,[67] from Dari in Afghanistan[68] to Japan.[69] Some studies have been done devoted to the spread of proverbs in certain regions, such as India and her neighbors[70] and Europe.[71] An extreme example of the borrowing and spread of proverbs was the work done to create a corpus of proverbs for Esperanto, where all the proverbs were translated from other languages.[72]
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It is often not possible to trace the direction of borrowing a proverb between languages. This is complicated by the fact that the borrowing may have been through plural languages. In some cases, it is possible to make a strong case for discerning the direction of the borrowing based on an artistic form of the proverb in one language, but a prosaic form in another language. For example, in Ethiopia there is a proverb "Of mothers and water, there is none evil." It is found in Amharic, Alaaba language, and Oromo, three languages of Ethiopia:
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The Oromo version uses poetic features, such as the initial ha in both clauses with the final -aa in the same word, and both clauses ending with -an. Also, both clauses are built with the vowel a in the first and last words, but the vowel i in the one syllable central word. In contrast, the Amharic and Alaaba versions of the proverb show little evidence of sound-based art.
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However, not all languages have proverbs. Proverbs are (nearly) universal across Europe, Asia, and Africa. Some languages in the Pacific have them, such as Maori.[74] Other Pacific languages do not, e.g. "there are no proverbs in Kilivila" of the Trobriand Islands.[75] However, in the New World, there are almost no proverbs: "While proverbs abound in the thousands in most cultures of the world, it remains a riddle why the Native Americans have hardly any proverb tradition at all."[76] Hakamies has examined the matter of whether proverbs are found universally, a universal genre, concluding that they are not.[77]
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Proverbs are used in conversation by adults more than children, partially because adults have learned more proverbs than children. Also, using proverbs well is a skill that is developed over years. Additionally, children have not mastered the patterns of metaphorical expression that are invoked in proverb use. Proverbs, because they are indirect, allow a speaker to disagree or give advice in a way that may be less offensive. Studying actual proverb use in conversation, however, is difficult since the researcher must wait for proverbs to happen.[78] An Ethiopian researcher, Tadesse Jaleta Jirata, made headway in such research by attending and taking notes at events where he knew proverbs were expected to be part of the conversations.[79]
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Many authors have used proverbs in their writings, for a very wide variety of literary genres: epics,[80][81][82][83] novels,[84][85] poems,[86] short stories.[87]
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Probably the most famous user of proverbs in novels is J. R. R. Tolkien in his The Hobbit and The Lord of the Rings series.[24][25][88][89] Herman Melville is noted for creating proverbs in Moby Dick[90] and in his poetry.[91][92] Also, C. S. Lewis created a dozen proverbs in The Horse and His Boy,[93] and Mercedes Lackey created dozens for her invented Shin'a'in and Tale'edras cultures;[94] Lackey's proverbs are notable in that they are reminiscent to those of Ancient Asia - e.g. "Just because you feel certain an enemy is lurking behind every bush, it doesn't follow that you are wrong" is like to "Before telling secrets on the road, look in the bushes." These authors are notable for not only using proverbs as integral to the development of the characters and the story line, but also for creating proverbs.[93]
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Among medieval literary texts, Geoffrey Chaucer's Troilus and Criseyde plays a special role because Chaucer's usage seems to challenge the truth value of proverbs by exposing their epistemological unreliability.[95] Rabelais used proverbs to write an entire chapter of Gargantua.[96]
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The patterns of using proverbs in literature can change over time. A study of "classical Chinese novels" found proverb use as frequently as one proverb every 3,500 words in Water Margin (Sui-hu chuan) and one proverb every 4,000 words in Wen Jou-hsiang. But modern Chinese novels have fewer proverbs by far.[97]
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Proverbs (or portions of them) have been the inspiration for titles of books: The Bigger they Come by Erle Stanley Gardner, and Birds of a Feather (several books with this title), Devil in the Details (multiple books with this title). Sometimes a title alludes to a proverb, but does not actually quote much of it, such as The Gift Horse's Mouth by Robert Campbell. Some books or stories have titles that are twisted proverbs, anti-proverbs, such as No use dying over spilled milk,[98] When life gives you lululemons,[99] and two books titled Blessed are the Cheesemakers.[100] The twisted proverb of last title was also used in the Monty Python movie Life of Brian, where a person mishears one of Jesus Christ's beatitudes, "I think it was 'Blessed are the cheesemakers.'"
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Some books and stories are built around a proverb. Some of Tolkien's books have been analyzed as having "governing proverbs" where "the acton of a book turns on or fulfills a proverbial saying."[101] Some stories have been written with a proverb overtly as an opening, such as "A stitch in time saves nine" at the beginning of "Kitty's Class Day", one of Louisa May Alcott's Proverb Stories. Other times, a proverb appears at the end of a story, summing up a moral to the story, frequently found in Aesop's Fables, such as "Heaven helps those who help themselves" from Hercules and the Wagoner.[102] In a novel by the Ivorian novelist Ahmadou Kourouma, "proverbs are used to conclude each chapter".[103]
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Proverbs have also been used strategically by poets.[104] Sometimes proverbs (or portions of them or anti-proverbs) are used for titles, such as "A bird in the bush" by Lord Kennet and his stepson Peter Scott and "The blind leading the blind" by Lisa Mueller. Sometimes, multiple proverbs are important parts of poems, such as Paul Muldoon's "Symposium", which begins "You can lead a horse to water but you can't make it hold its nose to the grindstone and hunt with the hounds. Every dog has a stitch in time..." In Finnish there are proverb poems written hundreds of years ago.[105] The Turkish poet Refiki wrote an entire poem by stringing proverbs together, which has been translated into English poetically yielding such verses as "Be watchful and be wary, / But seldom grant a boon; / The man who calls the piper / Will also call the tune."[106] Eliza Griswold also created a poem by stringing proverbs together, Libyan proverbs translated into English.[107]
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Because proverbs are familiar and often pointed, they have been used by a number of hip-hop poets. This has been true not only in the USA, birthplace of hip-hop, but also in Nigeria. Since Nigeria is so multilingual, hip-hop poets there use proverbs from various languages, mixing them in as it fits their need, sometimes translating the original. For example,
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"They forget say ogbon ju agbaralo
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They forget that wisdom is greater than power"[108]
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Some authors have bent and twisted proverbs, creating anti-proverbs, for a variety of literary effects. For example, in the Harry Potter novels, J. K. Rowling reshapes a standard English proverb into "It's no good crying over spilt potion" and Dumbledore advises Harry not to "count your owls before they are delivered".[109] In a slightly different use of reshaping proverbs, in the Aubrey–Maturin series of historical naval novels by Patrick O'Brian, Capt. Jack Aubrey humorously mangles and mis-splices proverbs, such as "Never count the bear's skin before it is hatched" and "There's a good deal to be said for making hay while the iron is hot."[110] Earlier than O'Brian's Aubrey, Beatrice Grimshaw also used repeated splicings of proverbs in the mouth of an eccentric marquis to create a memorable character in The Sorcerer's Stone, such as "The proof of the pudding sweeps clean" (p. 109) and "A stitch in time is as good as a mile" (p. 97).[111]
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Because proverbs are so much a part of the language and culture, authors have sometimes used proverbs in historical fiction effectively, but anachronistically, before the proverb was actually known. For example, the novel Ramage and the Rebels, by Dudley Pope is set in approximately 1800. Captain Ramage reminds his adversary "You are supposed to know that it is dangerous to change horses in midstream" (p. 259), with another allusion to the same proverb three pages later. However, the proverb about changing horses in midstream is reliably dated to 1864, so the proverb could not have been known or used by a character from that period.[112]
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Some authors have used so many proverbs that there have been entire books written cataloging their proverb usage, such as Charles Dickens,[113] Agatha Christie,[114] George Bernard Shaw,[115] Miguel de Cervantes,[116][117] and Friedrich Nietzsche.[118]
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On the non-fiction side, proverbs have also been used by authors for articles that have no connection to the study of proverbs. Some have been used as the basis for book titles, e.g. I Shop, Therefore I Am: Compulsive Buying and the Search for Self by April Lane Benson. Some proverbs been used as the basis for article titles, though often in altered form: "All our eggs in a broken basket: How the Human Terrain System is undermining sustainable military cultural competence"[119] and "Should Rolling Stones Worry About Gathering Moss?",[120] "Between a Rock and a Soft Place",[121] and the pair "Verbs of a feather flock together"[122] and "Verbs of a feather flock together II".[123] Proverbs have been noted as common in subtitles of articles[124] such as "Discontinued intergenerational transmission of Czech in Texas: 'Hindsight is better than foresight'."[125] Also, the reverse is found with a proverb (complete or partial) as the title, then an explanatory subtitle, "To Change or Not to Change Horses: The World War II Elections".[126] Many authors have cited proverbs as epigrams at the beginning of their articles, e.g. "'If you want to dismantle a hedge, remove one thorn at a time' Somali proverb" in an article on peacemaking in Somalia.[127] An article about research among the Māori used a Māori proverb as a title, then began the article with the Māori form of the proverb as an epigram "Set the overgrown bush alight and the new flax shoots will spring up", followed by three paragraphs about how the proverb served as a metaphor for the research and the present context.[128] A British proverb has even been used as the title for a doctoral dissertation: Where there is muck there is brass.[129]
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Similarly to other forms of literature, proverbs have also been used as important units of language in drama and films. This is true from the days of classical Greek works[130] to old French[131] to Shakespeare,[132] to 19th Century Spanish,[133] to today. The use of proverbs in drama and film today is still found in languages around the world, including Yorùbá [134] and Igbo[135] of Nigeria.
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A film that makes rich use of proverbs is Forrest Gump, known for both using and creating proverbs.[136][137] Other studies of the use of proverbs in film include work by Kevin McKenna on the Russian film Aleksandr Nevsky,[138] Haase's study of an adaptation of Little Red Riding Hood,[139] Elias Dominguez Barajas on the film Viva Zapata!,[140] and Aboneh Ashagrie on The Athlete (a movie in Amharic about Abebe Bikila).[141]
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Television programs have also been named with reference to proverbs, usually shortened, such Birds of a Feather and Diff'rent Strokes.
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In the case of Forrest Gump, the screenplay by Eric Roth had more proverbs than the novel by Winston Groom, but for The Harder They Come, the reverse is true, where the novel derived from the movie by Michael Thelwell has many more proverbs than the movie.[142]
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Éric Rohmer, the French film director, directed a series of films, the "Comedies and Proverbs", where each film was based on a proverb: The Aviator's Wife, The Perfect Marriage, Pauline at the Beach, Full Moon in Paris (the film's proverb was invented by Rohmer himself: "The one who has two wives loses his soul, the one who has two houses loses his mind."), The Green Ray, Boyfriends and Girlfriends.[143]
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Movie titles based on proverbs include Murder Will Out (1939 film), Try, Try Again, and The Harder They Fall. A twisted anti-proverb was the title for a Three Stooges film, A Bird in the Head. The title of an award-winning Turkish film, Three Monkeys, also invokes a proverb, though the title does not fully quote it.
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They have also been used as the titles of plays:[144] Baby with the Bathwater by Christopher Durang, Dog Eat Dog by Mary Gallagher, and The Dog in the Manger by Charles Hale Hoyt. The use of proverbs as titles for plays is not, of course, limited to English plays: Il faut qu'une porte soit ouverte ou fermée (A door must be open or closed) by Paul de Musset. Proverbs have also been used in musical dramas, such as The Full Monty, which has been shown to use proverbs in clever ways.[145] In the lyrics for Beauty and the Beast, Gaston plays with three proverbs in sequence, "All roads lead to.../The best things in life are.../All's well that ends with...me."
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Proverbs are often poetic in and of themselves, making them ideally suited for adapting into songs. Proverbs have been used in music from opera to country to hip-hop. Proverbs have also been used in music in many languages, such as the Akan language[146] the Igede language,[147] and Spanish.[148]
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In English the proverb (or rather the beginning of the proverb), If the shoe fits has been used as a title for three albums and five songs. Other English examples of using proverbs in music[149] include Elvis Presley's Easy come, easy go, Harold Robe's Never swap horses when you're crossing a stream, Arthur Gillespie's Absence makes the heart grow fonder, Bob Dylan's Like a rolling stone, Cher's Apples don't fall far from the tree. Lynn Anderson made famous a song full of proverbs, I never promised you a rose garden (written by Joe South). In choral music, we find Michael Torke's Proverbs for female voice and ensemble. A number of Blues musicians have also used proverbs extensively.[150][151] The frequent use of proverbs in Country music has led to published studies of proverbs in this genre.[152][153] The Reggae artist Jahdan Blakkamoore has recorded a piece titled Proverbs Remix. The opera Maldobrìe contains careful use of proverbs.[154] An extreme example of many proverbs used in composing songs is a song consisting almost entirely of proverbs performed by Bruce Springsteen, "My best was never good enough".[155] The Mighty Diamonds recorded a song called simply "Proverbs".[156]
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The band Fleet Foxes used the proverb painting Netherlandish Proverbs for the cover of their eponymous album Fleet Foxes.[157]
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In addition to proverbs being used in songs themselves, some rock bands have used parts of proverbs as their names, such as the Rolling Stones, Bad Company, The Mothers of Invention, Feast or Famine, Of Mice and Men. There have been at least two groups that called themselves "The Proverbs", and there is a hip-hop performer in South Africa known as "Proverb". In addition, many albums have been named with allusions to proverbs, such as Spilt milk (a title used by Jellyfish and also Kristina Train), The more things change by Machine Head, Silk purse by Linda Ronstadt, Another day, another dollar by DJ Scream Roccett, The blind leading the naked by Violent Femmes, What's good for the goose is good for the gander by Bobby Rush, Resistance is Futile by Steve Coleman, Murder will out by Fan the Fury. The proverb Feast or famine has been used as an album title by Chuck Ragan, Reef the Lost Cauze, Indiginus, and DaVinci. Whitehorse mixed two proverbs for the name of their album Leave no bridge unburned. The band Splinter Group released an album titled When in Rome, Eat Lions. The band Downcount used a proverb for the name of their tour, Come and take it.[158]
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From ancient times, people around the world have recorded proverbs in visual form. This has been done in two ways. First, proverbs have been written to be displayed, often in a decorative manner, such as on pottery, cross-stitch, murals,[159][160] kangas (East African women's wraps),[161] quilts,[162] a stained glass window,[63] and graffiti.[163]
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Secondly, proverbs have often been visually depicted in a variety of media, including paintings, etchings, and sculpture. Jakob Jordaens painted a plaque with a proverb about drunkenness above a drunk man wearing a crown, titled The King Drinks. Probably the most famous examples of depicting proverbs are the different versions of the paintings Netherlandish Proverbs by the father and son Pieter Bruegel the Elder and Pieter Brueghel the Younger, the proverbial meanings of these paintings being the subject of a 2004 conference, which led to a published volume of studies (Mieder 2004a). The same father and son also painted versions of The Blind Leading the Blind, a Biblical proverb. These and similar paintings inspired another famous painting depicting some proverbs and also idioms (leading to a series of additional paintings), such as Proverbidioms by T. E. Breitenbach. Another painting inspired by Bruegel's work is by the Chinese artist, Ah To, who created a painting illustrating 81 Cantonese sayings.[164] Corey Barksdale has produced a book of paintings with specific proverbs and pithy quotations.[165][self-published source?] The British artist Chris Gollon has painted a major work entitled "Big Fish Eat Little Fish", a title echoing Bruegel's painting Big Fishes Eat Little Fishes.[166]
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Sometimes well-known proverbs are pictured on objects, without a text actually quoting the proverb, such as the three wise monkeys who remind us "Hear no evil, see no evil, speak no evil". When the proverb is well known, viewers are able to recognize the proverb and understand the image appropriately, but if viewers do not recognize the proverb, much of the effect of the image is lost. For example, there is a Japanese painting in the Bonsai museum in Saitama city that depicted flowers on a dead tree, but only when the curator learned the ancient (and no longer current) proverb "Flowers on a dead tree" did the curator understand the deeper meaning of the painting.[167]
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A study of school students found that students remembered proverbs better when there were visual representations of proverbs along with the verbal form of the proverbs.[168]
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A bibliography on proverbs in visual form has been prepared by Mieder and Sobieski (1999). Interpreting visual images of proverbs is subjective, but familiarity with the depicted proverb helps.[169]
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Some artists have used proverbs and anti-proverbs for titles of their paintings, alluding to a proverb rather than picturing it. For example, Vivienne LeWitt painted a piece titled "If the shoe doesn't fit, must we change the foot?", which shows neither foot nor shoe, but a woman counting her money as she contemplates different options when buying vegetables.[170]
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In 2018, 13 sculptures depicting Maltese proverbs were installed in open spaces of downtown Valletta.[171][172][173]
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Cartoonists, both editorial and pure humorists, have often used proverbs, sometimes primarily building on the text, sometimes primarily on the situation visually, the best cartoons combining both. Not surprisingly, cartoonists often twist proverbs, such as visually depicting a proverb literally or twisting the text as an anti-proverb.[174] An example with all of these traits is a cartoon showing a waitress delivering two plates with worms on them, telling the customers, "Two early bird specials... here ya go."[175]
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The traditional Three wise monkeys were depicted in Bizarro with different labels. Instead of the negative imperatives, the one with ears covered bore the sign "See and speak evil", the one with eyes covered bore the sign "See and hear evil", etc. The caption at the bottom read "The power of positive thinking."[176] Another cartoon showed a customer in a pharmacy telling a pharmacist, "I'll have an ounce of prevention."[177] The comic strip The Argyle Sweater showed an Egyptian archeologist loading a mummy on the roof of a vehicle, refusing the offer of a rope to tie it on, with the caption "A fool and his mummy are soon parted."[178] The comic One Big Happy showed a conversation where one person repeatedly posed part of various proverb and the other tried to complete each one, resulting in such humorous results as "Don't change horses... unless you can lift those heavy diapers."[179]
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Editorial cartoons can use proverbs to make their points with extra force as they can invoke the wisdom of society, not just the opinion of the editors.[180] In an example that invoked a proverb only visually, when a US government agency (GSA) was caught spending money extravagantly, a cartoon showed a black pot labeled "Congress" telling a black kettle labeled "GSA", "Stop wasting the taxpayers' money!"[181] It may have taken some readers a moment of pondering to understand it, but the impact of the message was the stronger for it.
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Cartoons with proverbs are so common that Wolfgang Mieder has published a collected volume of them, many of them editorial cartoons. For example, a German editorial cartoon linked a current politician to the Nazis, showing him with a bottle of swastika-labeled wine and the caption "In vino veritas".[182]
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One cartoonist very self-consciously drew and wrote cartoons based on proverbs for the University of Vermont student newspaper The Water Tower, under the title "Proverb place".[183]
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Proverbs are frequently used in advertising, often in slightly modified form.[184][185][186]
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Ford once advertised its Thunderbird with, "One drive is worth a thousand words" (Mieder 2004b: 84). This is doubly interesting since the underlying proverb behind this, "One picture is worth a thousand words," was originally introduced into the English proverb repertoire in an ad for televisions (Mieder 2004b: 83).
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A few of the many proverbs adapted and used in advertising include:
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The GEICO company has created a series of television ads that are built around proverbs, such as "A bird in the hand is worth two in the bush",[187] and "The pen is mightier than the sword",[188] "Pigs may fly/When pigs fly",[189] "If a tree falls in the forest...",[190] and "Words can never hurt you".[191] Doritos made a commercial based on the proverb, "When pigs fly."[192] Many advertisements that use proverbs shorten or amend them, such as, "Think outside the shoebox."
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Use of proverbs in advertising is not limited to the English language. Seda Başer Çoban has studied the use of proverbs in Turkish advertising.[193] Tatira has given a number of examples of proverbs used in advertising in Zimbabwe.[194] However, unlike the examples given above in English, all of which are anti-proverbs, Tatira's examples are standard proverbs. Where the English proverbs above are meant to make a potential customer smile, in one of the Zimbabwean examples "both the content of the proverb and the fact that it is phrased as a proverb secure the idea of a secure time-honored relationship between the company and the individuals". When newer buses were imported, owners of older buses compensated by painting a traditional proverb on the sides of their buses, "Going fast does not assure safe arrival".[195]
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There are often proverbs that contradict each other, such as "Look before you leap" and "He who hesitates is lost", or "Many hands make light work" and "Too many cooks spoil the broth". These have been labeled "counter proverbs"[196] or "antonymous proverbs".[197] When there are such counter proverbs, each can be used in its own appropriate situation, and neither is intended to be a universal truth.[198][199] Some pairs of proverbs are fully contradictory: “A messy desk is a sign of intelligence” and “A neat desk is a sign of a sick mind”.[22]
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The concept of "counter proverb" is more about pairs of contradictory proverbs than about the use of proverbs to counter each other in an argument. For example, from the Tafi language of Ghana, the following pair of proverbs are counter to each other but are each used in appropriate contexts, "A co-wife who is too powerful for you, you address her as your mother" and "Do not call your mother's co-wife your mother..."[200] In Nepali, there is a set of totally contradictory proverbs: "Religion is victorious and sin erodes" and "Religion erodes and sin is victorious".[201]
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Also, the following pair are counter proverbs from the Kasena of Ghana: "It is the patient person who will milk a barren cow" and "The person who would milk a barren cow must prepare for a kick on the forehead".[202] The two contradict each other, whether they are used in an argument or not (though indeed they were used in an argument). But the same work contains an appendix with many examples of proverbs used in arguing for contrary positions, but proverbs that are not inherently contradictory,[203] such as "One is better off with hope of a cow's return than news of its death" countered by "If you don't know a goat [before its death] you mock at its skin". Though this pair was used in a contradictory way in a conversation, they are not a set of "counter proverbs".[198]
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Discussing counter proverbs in the Badaga language, Hockings explained that in his large collection "a few proverbs are mutually contradictory... we can be sure that the Badagas do not see the matter that way, and would explain such apparent contradictions by reasoning that proverb x is used in one context, while y is used in quite another."[204] Comparing Korean proverbs, "when you compare two proverbs, often they will be contradictory." They are used for "a particular situation".[205]
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"Counter proverbs" are not the same as a "paradoxical proverb", a proverb that contains a seeming paradox.[206]
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In many cultures, proverbs are so important and so prominent that there are proverbs about proverbs, that is, "metaproverbs". The most famous one is from Yoruba of Nigeria, "Proverbs are the horses of speech, if communication is lost we use proverbs to find it," used by Wole Soyinka in Death and the King's Horsemen. In Mieder's bibliography of proverb studies, there are twelve publications listed as describing metaproverbs.[207] Other metaproverbs include:
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There is a growing interest in deliberately using proverbs to achieve goals, usually to support and promote changes in society. Proverbs have also been used for public health promotion, such as promoting breast feeding with a shawl bearing a Swahili proverb "Mother's milk is sweet".[227] Proverbs have also been applied for helping people manage diabetes,[228] to combat prostitution,[229] and for community development.,[230] to resolve conflicts,[231][232] and to slow the transmission of HIV.[233]
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The most active field deliberately using proverbs is Christian ministry, where Joseph G. Healey and others have deliberately worked to catalyze the collection of proverbs from smaller languages and the application of them in a wide variety of church-related ministries, resulting in publications of collections[234] and applications.[235][236] This attention to proverbs by those in Christian ministries is not new, many pioneering proverb collections having been collected and published by Christian workers.[237][238][239][240]
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U.S. Navy Captain Edward Zellem pioneered the use of Afghan proverbs as a positive relationship-building tool during the war in Afghanistan, and in 2012 he published two bilingual collections[241][242] of Afghan proverbs in Dari and English, part of an effort of nationbuilding, followed by a volume of Pashto proverbs in 2014.[243]
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There is a longstanding debate among proverb scholars as to whether the cultural values of specific language communities are reflected (to varying degree) in their proverbs. Many claim that the proverbs of a particular culture reflect the values of that specific culture, at least to some degree. Many writers have asserted that the proverbs of their cultures reflect their culture and values; this can be seen in such titles as the following: An introduction to Kasena society and culture through their proverbs,[244] Prejudice, power, and poverty in Haiti: a study of a nation's culture as seen through its proverbs,[245] Proverbiality and worldview in Maltese and Arabic proverbs,[246] Fatalistic traits in Finnish proverbs,[247] Vietnamese cultural patterns and values as expressed in proverbs,[248] The Wisdom and Philosophy of the Gikuyu proverbs: The Kihooto worldview,[249] Spanish Grammar and Culture through Proverbs,[250] and "How Russian Proverbs Present the Russian National Character".[251] Kohistani has written a thesis to show how understanding Afghan Dari proverbs will help Europeans understand Afghan culture.[252]
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However, a number of scholars argue that such claims are not valid. They have used a variety of arguments. Grauberg argues that since many proverbs are so widely circulated they are reflections of broad human experience, not any one culture's unique viewpoint.[253] Related to this line of argument, from a collection of 199 American proverbs, Jente showed that only 10 were coined in the USA, so that most of these proverbs would not reflect uniquely American values.[254] Giving another line of reasoning that proverbs should not be trusted as a simplistic guide to cultural values, Mieder once observed "proverbs come and go, that is, antiquated proverbs with messages and images we no longer relate to are dropped from our proverb repertoire, while new proverbs are created to reflect the mores and values of our time",[255] so old proverbs still in circulation might reflect past values of a culture more than its current values. Also, within any language's proverb repertoire, there may be "counter proverbs", proverbs that contradict each other on the surface[196] (see section above). When examining such counter proverbs, it is difficult to discern an underlying cultural value. With so many barriers to a simple calculation of values directly from proverbs, some feel "one cannot draw conclusions about values of speakers simply from the texts of proverbs".[256]
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Many outsiders have studied proverbs to discern and understand cultural values and world view of cultural communities.[257] These outsider scholars are confident that they have gained insights into the local cultures by studying proverbs, but this is not universally accepted.[254][258][259][260]
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[261][262]
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Seeking empirical evidence to evaluate the question of whether proverbs reflect a culture's values, some have counted the proverbs that support various values. For example, Moon lists what he sees as the top ten core cultural values of the Builsa society of Ghana, as exemplified by proverbs. He found that 18% of the proverbs he analyzed supported the value of being a member of the community, rather than being independent.[263] This was corroboration to other evidence that collective community membership is an important value among the Builsa. In studying Tajik proverbs, Bell notes that the proverbs in his corpus "Consistently illustrate Tajik values" and "The most often observed proverbs reflect the focal and specific values" discerned in the thesis.[264]
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A study of English proverbs created since 1900 showed in the 1960s a sudden and significant increase in proverbs that reflected more casual attitudes toward sex.[265] Since the 1960s was also the decade of the Sexual revolution, this shows a strong statistical link between the changed values of the decades and a change in the proverbs coined and used. Another study mining the same volume counted Anglo-American proverbs about religion to show that proverbs indicate attitudes toward religion are going downhill.[266]
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There are many examples where cultural values have been explained and illustrated by proverbs. For example, from India, the concept that birth determines one's nature "is illustrated in the oft-repeated proverb: there can be no friendship between grass-eaters and meat-eaters, between a food and its eater".[267] Proverbs have been used to explain and illustrate the Fulani cultural value of pulaaku.[268] But using proverbs to illustrate a cultural value is not the same as using a collection of proverbs to discern cultural values. In a comparative study between Spanish and Jordanian proverbs it is defined the social imagination for the mother as an archetype in the context of role transformation and in contrast with the roles of husband, son and brother, in two societies which might be occasionally associated with sexist and /or rural ideologies.[269]
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Some scholars have adopted a cautious approach, acknowledging at least a genuine, though limited, link between cultural values and proverbs: "The cultural portrait painted by proverbs may be fragmented, contradictory, or otherwise at variance with reality... but must be regarded not as accurate renderings but rather as tantalizing shadows of the culture which spawned them."[270] There is not yet agreement on the issue of whether, and how much, cultural values are reflected in a culture's proverbs.
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It is clear that the Soviet Union believed that proverbs had a direct link to the values of a culture, as they used them to try to create changes in the values of cultures within their sphere of domination. Sometimes they took old Russian proverbs and altered them into socialist forms.[271] These new proverbs promoted Socialism and its attendant values, such as atheism and collectivism, e.g. "Bread is given to us not by Christ, but by machines and collective farms" and "A good harvest is had only by a collective farm." They did not limit their efforts to Russian, but also produced "newly coined proverbs that conformed to socialist thought" in Tajik and other languages of the USSR.[272]
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Many proverbs from around the world address matters of ethics and expected of behavior. Therefore, it is not surprising that proverbs are often important texts in religions. The most obvious example is the Book of Proverbs in the Bible. Additional proverbs have also been coined to support religious values, such as the following from Dari of Afghanistan:[273] "In childhood you're playful, In youth you're lustful, In old age you're feeble, So when will you before God be worshipful?"
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Clearly proverbs in religion are not limited to monotheists; among the Badagas of India (Sahivite Hindus), there is a traditional proverb "Catch hold of and join with the man who has placed sacred ash [on himself]."[274] Proverbs are widely associated with large religions that draw from sacred books, but they are also used for religious purposes among groups with their own traditional religions, such as the Guji Oromo.[79] The broadest comparative study of proverbs across religions is The eleven religions and their proverbial lore, a comparative study. A reference book to the eleven surviving major religions of the world by Selwyn Gurney Champion, from 1945. Some sayings from sacred books also become proverbs, even if they were not obviously proverbs in the original passage of the sacred book.[275] For example, many quote "Be sure your sin will find you out" as a proverb from the Bible, but there is no evidence it was proverbial in its original usage (Numbers 32:23).
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Not all religious references in proverbs are positive, some are cynical, such as the Tajik, "Do as the mullah says, not as he does."[276] Also, note the Italian proverb, "One barrel of wine can work more miracles than a church full of saints". An Indian proverb is cynical about devotees of Hinduism, "[Only] When in distress, a man calls on Rama".[277] In the context of Tibetan Buddhism, some Ladakhi proverbs mock the lamas, e.g. "If the lama's own head does not come out cleanly, how will he do the drawing upwards of the dead?... used for deriding the immoral life of the lamas."[278] Proverbs do not have to explicitly mention religion or religious figures to be used to mock a religion, seen in the fact that in a collection of 555 proverbs from the Lur, a Muslim group in Iran, the explanations for 15 of them use illustrations that mock Muslim clerics.[279]
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Dammann wrote, "In the [African] traditional religions, specific religious ideas recede into the background... The influence of Islam manifests itself in African proverbs... Christian influences, on the contrary, are rare."[280] If widely true in Africa, this is likely due to the longer presence of Islam in many parts of Africa. Reflection of Christian values is common in Amharic proverbs of Ethiopia, an area that has had a presence of Christianity for well over 1,000 years. The Islamic proverbial reproduction may also be shown in the image of some animals such as the dog. Although dog is portrayed in many European proverbs as the most faithful friend of man, it is represented in some Islamic countries as impure, dirty, vile, cowardly, ungrateful and treacherous, in addition to links to negative human superstitions such as loneliness, indifference and bad luck.[281]
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Though much proverb scholarship is done by literary scholars, those studying the human mind have used proverbs in a variety of studies.[282] One of the earliest studies in this field is the Proverbs Test by Gorham, developed in 1956. A similar test is being prepared in German.[283] Proverbs have been used to evaluate dementia,[284][285][286] study the cognitive development of children,[287] measure the results of brain injuries,[288] and study how the mind processes figurative language.[46][289][290]
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The study of proverbs is called paremiology which has a variety of uses in the study of such topics as philosophy, linguistics, and folklore. There are several types and styles of proverbs which are analyzed within Paremiology as is the use and misuse of familiar expressions which are not strictly 'proverbial' in the dictionary definition of being fixed sentences
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Grigorii Permjakov[291] developed the concept of the core set of proverbs that full members of society know, what he called the "paremiological minimum" (1979). For example, an adult American is expected to be familiar with "Birds of a feather flock together", part of the American paremiological minimum. However, an average adult American is not expected to know "Fair in the cradle, foul in the saddle", an old English proverb that is not part of the current American paremiological minimum. Thinking more widely than merely proverbs, Permjakov observed "every adult Russian language speaker (over 20 years of age) knows no fewer than 800 proverbs, proverbial expressions, popular literary quotations and other forms of cliches".[292] Studies of the paremiological minimum have been done for a limited number of languages, including Russian,[293] Hungarian,[294][295] Czech,[296] Somali,[297] Nepali,[298] Gujarati,[299] Spanish,[300] Esperanto,[301] Polish,[302] Ukrainian.[303] Two noted examples of attempts to establish a paremiological minimum in America are by Haas (2008) and Hirsch, Kett, and Trefil (1988), the latter more prescriptive than descriptive. There is not yet a recognized standard method for calculating the paremiological minimum, as seen by comparing the various efforts to establish the paremiological minimum in a number of languages.[304]
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A seminal work in the study of proverbs is Archer Taylor's The Proverb (1931), later republished by Wolfgang Mieder with Taylor's Index included (1985/1934). A good introduction to the study of proverbs is Mieder's 2004 volume, Proverbs: A Handbook. Mieder has also published a series of bibliography volumes on proverb research, as well as a large number of articles and other books in the field. Stan Nussbaum has edited a large collection on proverbs of Africa, published on a CD, including reprints of out-of-print collections, original collections, and works on analysis, bibliography, and application of proverbs to Christian ministry (1998).[305] Paczolay has compared proverbs across Europe and published a collection of similar proverbs in 55 languages (1997). Mieder edits an academic journal of proverb study, Proverbium (ISSN 0743-782X), many back issues of which are available online.[306] A volume containing articles on a wide variety of topics touching on proverbs was edited by Mieder and Alan Dundes (1994/1981). Paremia is a Spanish-language journal on proverbs, with articles available online.[307] There are also papers on proverbs published in conference proceedings volumes from the annual Interdisciplinary Colloquium on Proverbs[308] in Tavira, Portugal. Mieder has published a two-volume International Bibliography of Paremiology and Phraseology, with a topical, language, and author index.[309] Mieder has published a bibliography of collections of proverbs from around the world.[310] A broad introduction to proverb study, Introduction to Paremiology, edited by Hrisztalina Hrisztova-Gotthardt and Melita Aleksa Varga has been published in both hardcover and free open access, with articles by a dozen different authors.[311]
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Erasmus (1466–1536)
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Bishop Samuel Ajayi Crowther (c.1809-1891)
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Wolfgang Mieder
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Mineke Schipper
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Galit Hasan Rokem
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Dora Sakayan
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The study of proverbs has been built by a number of notable scholars and contributors. Earlier scholars were more concerned with collecting than analyzing. Desiderius Erasmus was a Latin scholar (1466 – 1536), whose collection of Latin proverbs, known as Adagia, spread Latin proverbs across Europe.[312] Juan de Mal Lara was a 16th century Spanish scholar, one of his books being 1568 Philosophia vulgar, the first part of which contains one thousand and one sayings. Hernán Núñez published a collection of Spanish proverbs (1555).
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In the 19th century, a growing number of scholars published collections of proverbs, such as Samuel Adalberg who published collections of Yiddish proverbs (1888 & 1890) and Polish proverbs (1889–1894). Samuel Ajayi Crowther, the Anglican bishop in Nigeria, published a collection of Yoruba proverbs (1852). Elias Lönnrot published a collection of Finnish proverbs (1842).
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From the 20th century onwards, proverb scholars were involved in not only collecting proverbs, but also analyzing and comparing proverbs. Alan Dundes was a 20th century American folklorist whose scholarly output on proverbs led Wolfgang Mieder to refer to him as a "pioneering paremiologist".[313] Matti Kuusi was a 20th century Finnish paremiologist, the creator of the Matti Kuusi international type system of proverbs.[314] With encouragement from Archer Taylor,[315] he founded the journal Proverbium: Bulletin d'Information sur les Recherches Parémiologiques, published from 1965 to 1975 by the Society for Finnish Literature, which was later restarted as Proverbium: International Yearbook of Proverb Scholarship. Archer Taylor was a 20th century American scholar, best known for his "magisterial"[316] book The Proverb.[317] Dimitrios Loukatos was a 20th century Greek proverb scholar, author of such works as Aetiological Tales of Modern Greek Proverbs.[318] Arvo Krikmann (1939 – 2017) was an Estonian proverb scholar, whom Wolfgang Mieder called "one of the leading paremiologists in the world"[319] and "master folklorist and paremiologist".[320] Elisabeth Piirainen was a German scholar with 50 proverb-related publications.[321]
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Current proverb scholars have continued the trend to be involved in analysis as well as collection of proverbs. Claude Buridant is a 20th century French scholar whose work has concentrated on Romance languages.[322] Galit Hasan-Rokem is an Israeli scholar, associate editor of Proverbium: The yearbook of international proverb scholarship, since 1984. She has written on proverbs in Jewish traditions.[323] Joseph G. Healey is an American Catholic missionary in Kenya who has led a movement to sponsor African proverb scholars to collect proverbs from their own language communities.[324] This led Wolfgang Mieder to dedicate the "International Bibliography of New and Reprinted Proverb Collections" section of Proverbium 32 to Healey.[325] Barbara Kirshenblatt-Gimblett is a scholar of Jewish history and folklore, including proverbs.[326] Wolfgang Mieder is a German-born proverb scholar who has worked his entire academic career in the USA. He is the editor of ‘’Proverbium’’ and the author of the two volume International Bibliography of Paremiology and Phraseology.[327] He has been honored by four festschrift publications.[328][329][330][331] He has also been recognized by biographical publications that focused on his scholarship.[332][333] Dora Sakayan is a scholar who has written about German and Armenian studies, including Armenian Proverbs: A Paremiological Study with an Anthology of 2,500 Armenian Folk Sayings Selected and Translated into English.[334] An extensive introduction addresses the language and structure,[335] as well as the origin of Armenian proverbs (international, borrowed and specifically Armenian proverbs). Mineke Schipper is a Dutch scholar, best known for her book Never Marry a Woman with Big Feet - Women in Proverbs from Around the World.[336] Edward Zellem is an American proverb scholar who has edited books of Afghan proverbs, developed a method of collecting proverbs via the Web.[337]
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Websites related to the study of proverbs, and some that list regional proverbs:
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1 |
+
West Flanders (Dutch: West-Vlaanderen [ˌʋɛst ˈflaːndərə(n)] (listen);[4] West Flemish: West Vloandern; French: (Province de) Flandre-Occidentale; German: Westflandern) is the westernmost province of the Flemish Region, in Belgium. It is the only coastal Belgian province, facing the North Sea to the northwest. It has land borders with the Dutch province of Zeeland to the northeast, the Flemish province of East Flanders to the east, the Walloon province of Hainaut in the southeast and the French department of Nord to the west. Its capital is Bruges (Brugge). Other important cities are Kortrijk in the south and Ostend (Oostende) on the coast, Roeselare and Ypres (Ieper). The province has an area of 3,197 km2 (1,234 sq mi) which is divided into eight administrative districts (arrondissementen) containing 64 municipalities.
|
2 |
+
As of January 2019, West Flanders has a population of 1,195,796.[5]
|
3 |
+
|
4 |
+
The North Sea coast of Belgium, an important tourism destination, lies in West Flanders. A tram line runs the length of the coast, from De Panne on the French border to Knokke-Heist on the Dutch border. West Flanders has two seaports, the Port of Zeebrugge and the Port of Ostend.
|
5 |
+
|
6 |
+
West Flanders consists of the North Sea coast, followed by a very flat polder landscape. Only in the south are some small hills, with the Kemmelberg (159 m (522 ft)) being the highest point in the province. The Leie and IJzer are the main rivers. West Flanders is both the only Belgian province that borders both France and The Netherlands; it is the only province not landlocked.
|
7 |
+
|
8 |
+
The Gross domestic product (GDP) of the province was 46.9 billion € in 2018. GDP per capita adjusted for purchasing power was 34,700 € or 115% of the EU27 average in the same year.[6]
|
9 |
+
|
10 |
+
In the north of the province, most industry is concentrated in and around the cities of Bruges and Oostende. Both cities also have important seaports: the port of Bruges-Zeebrugge and the port of Ostend.
|
11 |
+
|
12 |
+
The south is known for its textile industry, with companies such as Beaulieu and Libeco-Lagae. The region around Kortrijk is called the "Dallas of Belgium" for its entrepreneurship.
|
13 |
+
|
14 |
+
Tourism is also an important industry in West Flanders. Major touristic attractions include the Belgian coast, the historic center of Bruges, the Yser Tower in Diksmuide and Flanders Fields, the World War I battlefields around Ypres.
|
15 |
+
|
16 |
+
The region has a very distinct dialect of Dutch called West Flemish. This distinctive dialect is a part of the strong sense of identity of the inhabitants of the region.
|
17 |
+
|
18 |
+
The city of Bruges has two clubs playing at the highest level of football in Belgium (Belgian First Division A), namely Club Brugge and Cercle Brugge.
|
19 |
+
|
20 |
+
Other cities with teams playing at the highest level are Kortrijk (KV Kortrijk), Ostend (KV Oostende) and Waregem (Zulte Waregem).
|
21 |
+
|
22 |
+
The teams Knack Randstad Roeselare and Prefaxis Menen play in the highest volleyball league. Knack Randstad Roeselare also plays in the CEV Champions League, the European Championsleague.
|
23 |
+
|
24 |
+
BC Oostende is a basketball team which plays in the major basketball league of Belgium.
|
25 |
+
|
26 |
+
The Province of West Flanders is divided into 8 administrative arrondissements.
|
27 |
+
|
28 |
+
Bruges
|
29 |
+
|
30 |
+
Diksmuide
|
31 |
+
|
32 |
+
Kortrijk
|
33 |
+
|
34 |
+
Ostend
|
35 |
+
|
36 |
+
Roeselare
|
37 |
+
|
38 |
+
Tielt
|
39 |
+
|
40 |
+
Veurne
|
41 |
+
|
42 |
+
Ypres
|
43 |
+
|
44 |
+
The Province of West Flanders is divided into 4 judicial arrondissements.
|
45 |
+
|
46 |
+
Bruges
|
47 |
+
|
48 |
+
Kortrijk
|
49 |
+
|
50 |
+
Veurne
|
51 |
+
|
52 |
+
Ypres
|
53 |
+
|
54 |
+
Municipalities that have city status have (city) after their names.
|
55 |
+
|
56 |
+
1. Alveringem
|
57 |
+
2. Anzegem
|
58 |
+
3. Ardooie
|
59 |
+
4. Avelgem
|
60 |
+
5. Beernem
|
61 |
+
6. Blankenberge (city)
|
62 |
+
7. Bredene
|
63 |
+
8. Bruges (city)
|
64 |
+
9. Damme (city)
|
65 |
+
10. De Haan
|
66 |
+
11. De Panne
|
67 |
+
12. Deerlijk
|
68 |
+
13. Dentergem
|
69 |
+
14. Diksmuide (city)
|
70 |
+
15. Gistel (city)
|
71 |
+
16. Harelbeke (city)
|
72 |
+
17. Heuvelland
|
73 |
+
18. Hooglede
|
74 |
+
19. Houthulst
|
75 |
+
20. Ichtegem
|
76 |
+
21. Ypres (Ieper) (city)
|
77 |
+
22. Ingelmunster
|
78 |
+
23. Izegem (city)
|
79 |
+
24. Jabbeke
|
80 |
+
25. Knokke-Heist
|
81 |
+
26. Koekelare
|
82 |
+
27. Koksijde
|
83 |
+
28. Kortemark
|
84 |
+
29. Kortrijk (city)
|
85 |
+
30. Kuurne
|
86 |
+
31. Langemark-Poelkapelle
|
87 |
+
32. Ledegem
|
88 |
+
33. Lendelede
|
89 |
+
34. Lichtervelde
|
90 |
+
35. Lo-Reninge (city)
|
91 |
+
36. Menen (city)
|
92 |
+
37. Mesen (city)
|
93 |
+
38. Meulebeke
|
94 |
+
39. Middelkerke
|
95 |
+
40. Moorslede
|
96 |
+
41. Nieuwpoort (city)
|
97 |
+
42. Ostend (city)
|
98 |
+
43. Oostkamp
|
99 |
+
44. Oostrozebeke
|
100 |
+
45. Oudenburg (city)
|
101 |
+
46. Pittem
|
102 |
+
47. Poperinge (city)
|
103 |
+
48. Roeselare (city)
|
104 |
+
49. Ruiselede
|
105 |
+
50. Spiere-Helkijn
|
106 |
+
51. Staden
|
107 |
+
52. Tielt (city)
|
108 |
+
53. Torhout (city)
|
109 |
+
54. Veurne (city)
|
110 |
+
55. Vleteren
|
111 |
+
56. Waregem (city)
|
112 |
+
57. Wervik (city)
|
113 |
+
58. Wevelgem
|
114 |
+
59. Wielsbeke
|
115 |
+
60. Wingene
|
116 |
+
61. Zedelgem
|
117 |
+
62. Zonnebeke
|
118 |
+
63. Zuienkerke
|
119 |
+
64. Zwevegem
|
en/4842.html.txt
ADDED
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|
1 |
+
|
2 |
+
|
3 |
+
Quebec (/k(w)ɪˈbɛk/ (listen);[9] French: Québec [kebɛk] (listen))[10] is one of the thirteen provinces and territories of Canada. It is bordered to the west by the province of Ontario and the bodies of water James Bay and Hudson Bay; to the north by Hudson Strait and Ungava Bay; to the east by the Gulf of Saint Lawrence and the province of Newfoundland and Labrador; and to the south by the province of New Brunswick and the U.S. states of Maine, New Hampshire, Vermont, and New York. It also shares maritime borders with Nunavut, Prince Edward Island, and Nova Scotia. Quebec is Canada's largest province by area and its second-largest administrative division; only the territory of Nunavut is larger. It is historically and politically considered to be part of Central Canada (with Ontario).
|
4 |
+
|
5 |
+
Quebec is the second-most populous province of Canada, after Ontario.[11] It is the only one to have a predominantly French-speaking population, with French as the sole provincial official language. Most inhabitants live in urban areas near the Saint Lawrence River between Montreal and Quebec City, the capital. Approximately half of Quebec residents live in the Greater Montreal Area, including the Island of Montreal. English-speaking communities and English-language institutions are concentrated in the west of the island of Montreal but are also significantly present in the Outaouais, Eastern Townships, and Gaspé regions. The Nord-du-Québec region, occupying the northern half of the province, is sparsely populated and inhabited primarily by Aboriginal peoples.[12]
|
6 |
+
|
7 |
+
The climate around the major cities is four-seasons continental with cold and snowy winters combined with warm to hot humid summers, but farther north long winter seasons dominate and as a result the northern areas of the province are marked by tundra conditions.[13] Even in central Quebec, at comparatively southerly latitudes, winters are severe in inland areas.
|
8 |
+
|
9 |
+
Quebec independence debates have played a large role in the politics of the province. Parti Québécois governments held referendums on sovereignty in 1980 and 1995. Although neither passed, the 1995 referendum saw the highest voter turnout in Quebec history, at over 93%, and only failed by less than 1%.[14] In 2006, the House of Commons of Canada passed a symbolic motion recognizing the "Québécois as a nation within a united Canada".[15][16]
|
10 |
+
|
11 |
+
While the province's substantial natural resources have long been the mainstay of its economy, sectors of the knowledge economy such as aerospace, information and communication technologies, biotechnology, and the pharmaceutical industry also play leading roles. These many industries have all contributed to helping Quebec become an economically influential province within Canada, second only to Ontario in economic output.[17]
|
12 |
+
|
13 |
+
The name "Québec", which comes from the Algonquin or Ojibwe word kébec meaning "where the river narrows", originally referred to the area around Quebec City where the Saint Lawrence River narrows to a cliff-lined gap. Early variations in the spelling of the name included Québecq (Levasseur, 1601) and Kébec (Lescarbot, 1609).[18] French explorer Samuel de Champlain chose the name Québec in 1608 for the colonial outpost he would use as the administrative seat for the French colony of New France.[19] The province is sometimes referred to as "La belle province" ("The beautiful province").
|
14 |
+
|
15 |
+
The Province of Quebec was founded in the Royal Proclamation of 1763 after the Treaty of Paris formally transferred the French colony of Canada[20] to Britain after the Seven Years' War. The proclamation restricted the province to an area along the banks of the Saint Lawrence River. The Quebec Act of 1774 expanded the territory of the province to include the Great Lakes and the Ohio River Valley and south of Rupert's Land, more or less restoring the borders previously existing under French rule before the Conquest of 1760.[21] The Treaty of Paris (1783) ceded territories south of the Great Lakes to the United States.[22] After the Constitutional Act of 1791, the territory was divided between Lower Canada (present-day Quebec) and Upper Canada (present-day Ontario), with each being granted an elected legislative assembly.[23] In 1840, these become Canada East and Canada West after the British Parliament unified Upper and Lower Canada into the Province of Canada.[24] This territory was redivided into the Provinces of Quebec and Ontario at Confederation in 1867.[25] Each became one of the first four provinces.
|
16 |
+
|
17 |
+
In 1870, Canada purchased Rupert's Land from the Hudson's Bay Company and over the next few decades the Parliament of Canada transferred to Quebec portions of this territory that would more than triple the size of the province.[26] In 1898, the Canadian Parliament passed the first Quebec Boundary Extension Act that expanded the provincial boundaries northward to include the lands of the local aboriginal peoples.[27] This was followed by the addition of the District of Ungava through the Quebec Boundaries Extension Act of 1912 that added the northernmost lands of the Inuit to create the modern Province of Quebec.[27] In 1927, the border between Quebec and Newfoundland and Labrador was established by the British Judicial Committee of the Privy Council. Quebec officially disputes this boundary.[28]
|
18 |
+
|
19 |
+
Located in the eastern part of Canada, and (from a historical and political perspective) part of Central Canada, Quebec occupies a territory nearly three times the size of France or Texas, most of which is very sparsely populated.[29] Its topography is very different from one region to another due to the varying composition of the ground, the climate (latitude and altitude), and the proximity to water. The Saint Lawrence Lowland and the Appalachians are the two main topographic regions in southern Quebec, while the Canadian Shield occupies most of central and northern Quebec.[30]
|
20 |
+
|
21 |
+
Quebec has one of the world's largest reserves of fresh water,[31] occupying 12% of its surface.[32] It has 3% of the world's renewable fresh water, whereas it has only 0.1% of its population.[33] More than half a million lakes,[31] including 30 with an area greater than 250 square kilometres (97 sq mi), and 4,500 rivers[31] pour their torrents into the Atlantic Ocean, through the Gulf of Saint Lawrence and the Arctic Ocean, by James, Hudson, and Ungava bays. The largest inland body of water is the Caniapiscau Reservoir, created in the realization of the James Bay Project to produce hydroelectric power. Lake Mistassini is the largest natural lake in Quebec.[34]
|
22 |
+
|
23 |
+
The Saint Lawrence River has some of the world's largest sustaining inland Atlantic ports at Montreal (the province's largest city), Trois-Rivières, and Quebec City (the capital). Its access to the Atlantic Ocean and the interior of North America made it the base of early French exploration and settlement in the 17th and 18th centuries. Since 1959, the Saint Lawrence Seaway has provided a navigable link between the Atlantic Ocean and the Great Lakes. Northeast of Quebec City, the river broadens into the world's largest estuary, the feeding site of numerous species of whales, fish, and seabirds.[35] The river empties into the Gulf of Saint Lawrence. This marine environment sustains fisheries and smaller ports in the Lower Saint Lawrence (Bas-Saint-Laurent), Lower North Shore (Côte-Nord), and Gaspé (Gaspésie) regions of the province. The Saint Lawrence River with its estuary forms the basis of Quebec's development through the centuries. Other notable rivers include the Ashuapmushuan, Chaudière, Gatineau, Manicouagan, Ottawa, Richelieu, Rupert, Saguenay, Saint-François, and Saint-Maurice.
|
24 |
+
|
25 |
+
Quebec's highest point at 1,652 metres is Mont d'Iberville, known in English as Mount Caubvick, located on the border with Newfoundland and Labrador in the northeastern part of the province, in the Torngat Mountains.[36] The most populous physiographic region is the Saint Lawrence Lowland. It extends northeastward from the southwestern portion of the province along the shores of the Saint Lawrence River to the Quebec City region, limited to the North by the Laurentian Mountains and to the South by the Appalachians. It mainly covers the areas of the Centre-du-Québec, Laval, Montérégie and Montreal, the southern regions of the Capitale-Nationale, Lanaudière, Laurentides, Mauricie and includes Anticosti Island, the Mingan Archipelago,[37] and other small islands of the Gulf of St. Lawrence lowland forests ecoregion.[38] Its landscape is low-lying and flat, except for isolated igneous outcrops near Montreal called the Monteregian Hills, formerly covered by the waters of Lake Champlain. The Oka hills also rise from the plain. Geologically, the lowlands formed as a rift valley about 100 million years ago and are prone to infrequent but significant earthquakes.[30] The most recent layers of sedimentary rock were formed as the seabed of the ancient Champlain Sea at the end of the last ice age about 14,000 years ago.[39] The combination of rich and easily arable soils and Quebec's relatively warm climate makes this valley the most prolific agricultural area of Quebec province. Mixed forests provide most of Canada's springtime maple syrup crop. The rural part of the landscape is divided into narrow rectangular tracts of land that extend from the river and date back to settlement patterns in 17th century New France.
|
26 |
+
|
27 |
+
More than 95% of Quebec's territory lies within the Canadian Shield.[40] It is generally a quite flat and exposed mountainous terrain interspersed with higher points such as the Laurentian Mountains in southern Quebec, the Otish Mountains in central Quebec and the Torngat Mountains near Ungava Bay. The topography of the Shield has been shaped by glaciers from the successive ice ages, which explains the glacial deposits of boulders, gravel and sand, and by sea water and post-glacial lakes that left behind thick deposits of clay in parts of the Shield. The Canadian Shield also has a complex hydrological network of perhaps a million lakes, bogs, streams and rivers. It is rich in the forestry, mineral and hydro-electric resources that are a mainstay of the Quebec economy. Primary industries sustain small cities in regions of Abitibi-Témiscamingue, Saguenay–Lac-Saint-Jean, and Côte-Nord.
|
28 |
+
|
29 |
+
The Labrador Peninsula is covered by the Laurentian Plateau (or Canadian Shield), dotted with mountains such as Otish Mountains. The Ungava Peninsula is notably composed of D'Youville mountains, Puvirnituq mountains and Pingualuit crater. While low and medium altitude peak from western Quebec to the far north, high altitudes mountains emerge in the Capitale-Nationale region to the extreme east, along its longitude. In the Labrador Peninsula portion of the Shield, the far northern region of Nunavik includes the Ungava Peninsula and consists of flat Arctic tundra inhabited mostly by the Inuit. Further south lie the subarctic taiga of the Eastern Canadian Shield taiga ecoregion and the boreal forest of the Central Canadian Shield forests, where spruce, fir, and poplar trees provide raw materials for Quebec's pulp and paper and lumber industries. Although the area is inhabited principally by the Cree, Naskapi, and Innu First Nations, thousands of temporary workers reside at Radisson to service the massive James Bay Hydroelectric Project on the La Grande and Eastmain rivers. The southern portion of the shield extends to the Laurentians, a mountain range just north of the Saint Lawrence Lowland, that attracts local and international tourists to ski hills and lakeside resorts.
|
30 |
+
|
31 |
+
The Appalachian region of Quebec has a narrow strip of ancient mountains along the southeastern border of Quebec. The Appalachians are actually a huge chain that extends from Alabama to Newfoundland. In between, it covers in Quebec near 800 km (497 mi), from the Montérégie hills to the Gaspé Peninsula. In western Quebec, the average altitude is about 500 metres, while in the Gaspé Peninsula, the Appalachian peaks (especially the Chic-Choc) are among the highest in Quebec, exceeding 1000 metres.
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In general, the climate of Quebec is cold and humid.[41] The climate of the province is largely determined by its latitude, maritime and elevation influences.[41] According to the Köppen climate classification, Quebec has three main climate regions.[41] Southern and western Quebec, including most of the major population centres and areas south of 51oN, have a humid continental climate (Köppen climate classification Dfb) with four distinct seasons having warm to occasionally hot and humid summers and often very cold and snowy winters.[41][42] The main climatic influences are from western and northern Canada and move eastward, and from the southern and central United States that move northward. Because of the influence of both storm systems from the core of North America and the Atlantic Ocean, precipitation is abundant throughout the year, with most areas receiving more than 1,000 millimetres (39 in) of precipitation, including over 300 centimetres (120 in) of snow in many areas.[43] During the summer, severe weather patterns (such as tornadoes and severe thunderstorms) occur occasionally.[44] Most of central Quebec, ranging from 51 to 58 degrees North has a subarctic climate (Köppen Dfc).[41] Winters are long, very cold, and snowy, and among the coldest in eastern Canada, while summers are warm but very short due to the higher latitude and the greater influence of Arctic air masses. Precipitation is also somewhat less than farther south, except at some of the higher elevations. The northern regions of Quebec have an arctic climate (Köppen ET), with very cold winters and short, much cooler summers.[41] The primary influences in this region are the Arctic Ocean currents (such as the Labrador Current) and continental air masses from the High Arctic.
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The four calendar seasons in Quebec are spring, summer, autumn and winter, with conditions differing by region. They are then differentiated according to the insolation, temperature, and precipitation of snow and rain.[45]
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36 |
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At Quebec City, the length of the daily sunshine varies from 8:37 hrs in December to 15:50 hrs in June; the annual variation is much greater (from 4:54 to 19:29 hrs) at the northern tip of the province.[46] From temperate zones to the northern territories of the Far North, the brightness varies with latitude, as well as the Northern Lights and midnight sun.
|
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Quebec is divided into four climatic zones: arctic, subarctic, humid continental and East maritime. From south to north, average temperatures range in summer between 25 and 5 °C (77 and 41 °F) and, in winter, between −10 and −25 °C (14 and −13 °F).[47][48] In periods of intense heat and cold, temperatures can reach 35 °C (95 °F) in the summer[49] and −40 °C (−40 °F) during the Quebec winter,[49] They may vary depending on the Humidex or Wind chill. The all time record high was 40.0 °C (104.0 °F) and the all time record low was −51.0 °C (−59.8 °F).[50]
|
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The all-time record of the greatest precipitation in winter was established in winter 2007–2008, with more than five metres[51] of snow in the area of Quebec City, while the average amount received per winter is around three metres.[52] March 1971, however, saw the "Century's Snowstorm" with more than 40 centimetres (16 in) in Montreal to 80 centimetres (31 in) in Mont Apica of snow within 24 hours in many regions of southern Quebec. Also, the winter of 2010 was the warmest and driest recorded in more than 60 years.[53]
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The large land wildlife is mainly composed of the white-tailed deer, the moose, the muskox, the caribou, the American black bear and the polar bear. The average land wildlife includes the cougar, the coyote, the eastern wolf, the bobcat, the Arctic fox, the fox, etc. The small animals seen most commonly include the eastern grey squirrel, the snowshoe hare, the groundhog, the skunk, the raccoon, the chipmunk and the Canadian beaver.
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Biodiversity of the estuary and gulf of Saint Lawrence River[55] consists of an aquatic mammal wildlife, of which most goes upriver through the estuary and the Saguenay–St. Lawrence Marine Park until the Île d'Orléans (French for Orleans Island), such as the blue whale, the beluga, the minke whale and the harp seal (earless seal). Among the Nordic marine animals, there are two particularly important to cite: the walrus and the narwhal.[56]
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Inland waters are populated by small to large fresh water fish, such as the largemouth bass, the American pickerel, the walleye, the Acipenser oxyrinchus, the muskellunge, the Atlantic cod, the Arctic char, the brook trout, the Microgadus tomcod (tomcod), the Atlantic salmon, the rainbow trout, etc.[57]
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Among the birds commonly seen in the southern inhabited part of Quebec, there are the American robin, the house sparrow, the red-winged blackbird, the mallard, the common grackle, the blue jay, the American crow, the black-capped chickadee, some warblers and swallows, the starling and the rock pigeon, the latter two having been introduced in Quebec and are found mainly in urban areas.[58] Avian fauna includes birds of prey like the golden eagle, the peregrine falcon, the snowy owl and the bald eagle. Sea and semi-aquatic birds seen in Quebec are mostly the Canada goose, the double-crested cormorant, the northern gannet, the European herring gull, the great blue heron, the sandhill crane, the Atlantic puffin and the common loon.[59]
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Many more species of land, maritime or avian wildlife are seen in Quebec, but most of the Quebec-specific species and the most commonly seen species are listed above.
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Some livestock have the title of "Québec heritage breed", namely the Canadian horse, the Chantecler chicken and the Canadian cow.[60] Moreover, in addition to food certified as "organic", Charlevoix lamb is the first local Quebec product whose geographical indication is protected.[61] Livestock production also includes the pig breeds Landrace, Duroc and Yorkshire[62] and many breeds of sheep[63] and cattle.
|
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The Wildlife Foundation of Quebec and the Data Centre on Natural Heritage of Quebec (CDPNQ)(French acronym)[64] are the main agencies working with officers for wildlife conservation in Quebec.
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Given the geology of the province and its different climates, there is an established number of large areas of vegetation in Quebec. These areas, listed in order from the northernmost to the southernmost are: the tundra, the taiga, the Canadian boreal forest (coniferous), mixed forest and Deciduous forest.[40]
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On the edge of the Ungava Bay and Hudson Strait is the tundra, whose flora is limited to a low vegetation of lichen with only less than 50 growing days a year. The tundra vegetation survives an average annual temperature of −8 °C (18 °F). The tundra covers more than 24% of the area of Quebec.[40] Further south, the climate is conducive to the growth of the Canadian boreal forest, bounded on the north by the taiga.
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Not as arid as the tundra, the taiga is associated with the sub-Arctic regions of the Canadian Shield[65] and is characterized by a greater number of both plant (600) and animal (206) species, many of which live there all year. The taiga covers about 20% of the total area of Quebec.[40] The Canadian boreal forest is the northernmost and most abundant of the three forest areas in Quebec that straddle the Canadian Shield and the upper lowlands of the province. Given a warmer climate, the diversity of organisms is also higher, since there are about 850 plant species and 280 vertebrates species. The Canadian boreal forest covers 27% of the area of Quebec.[40] The mixed forest is a transition zone between the Canadian boreal forest and deciduous forest. By virtue of its transient nature, this area contains a diversity of habitats resulting in large numbers of plant (1000) and vertebrates (350) species, despite relatively cool temperatures. The ecozone mixed forest covers 11.5% of the area of Quebec and is characteristic of the Laurentians, the Appalachians and the eastern lowlands forests.[65] The third most northern forest area is characterized by deciduous forests. Because of its climate (average annual temperature of 7 °C (45 °F)), it is in this area that one finds the greatest diversity of species, including more than 1600 vascular plants and 440 vertebrates. Its relatively long growing season lasts almost 200 days and its fertile soils make it the centre of agricultural activity and therefore of urbanization of Quebec. Most of Quebec's population lives in this area of vegetation, almost entirely along the banks of the St. Lawrence. Deciduous forests cover approximately 6.6% of the area of Quebec.[40]
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The total forest area of Quebec is estimated at 750,300 square kilometres (289,700 sq mi).[66] From the Abitibi-Témiscamingue to the North Shore, the forest is composed primarily of conifers such as the Abies balsamea, the jack pine, the white spruce, the black spruce and the tamarack. Some species of deciduous trees such as the yellow birch appear when the river is approached in the south. The deciduous forest of the Saint Lawrence Lowlands is mostly composed of deciduous species such as the sugar maple, the red maple, the white ash, the American beech, the butternut (white walnut), the American elm, the basswood, the bitternut hickory and the northern red oak as well as some conifers such as the eastern white pine and the northern whitecedar. The distribution areas of the paper birch, the trembling aspen and the mountain ash cover more than half of Quebec territory.[67]
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At the time of first European contact and later colonization, Algonquian, Iroquois and Inuit nations controlled what is now Quebec.[68] Their lifestyles and cultures reflected the land on which they lived. Algonquians organized into seven political entities lived nomadic lives based on hunting, gathering, and fishing in the rugged terrain of the Canadian Shield (James Bay Cree, Innu, Algonquins) and Appalachian Mountains (Mi'kmaq, Abenaki).[69] St. Lawrence Iroquoians, a branch of the Iroquois, lived more settled lives, growing corn, beans and squash in the fertile soils of the St. Lawrence Valley. They appear to have been later supplanted by the Mohawk nation.[70] The Inuit continue to fish and hunt whale and seal in the harsh Arctic climate along the coasts of Hudson and Ungava Bay.[71] These people traded fur and food and sometimes warred with each other.
|
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Around 1522–1523, the Italian navigator Giovanni da Verrazzano persuaded King Francis I of France to commission an expedition to find a western route to Cathay (China). In 1534, Breton explorer Jacques Cartier planted a cross in the Gaspé Peninsula and claimed the land in the name of King Francis I.[73] It was the first province of New France. However, initial French attempts at settling the region met with failure.[73] French fishing fleets, however, continued to sail to the Atlantic coast and into the St. Lawrence River, making alliances with First Nations that would become important once France began to occupy the land.[74]
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+
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Samuel de Champlain was part of a 1603 expedition from France that travelled into the St. Lawrence River.[75] In 1608, he returned as head of an exploration party and founded Quebec City with the intention of making the area part of the French colonial empire.[76][77][78] Champlain's Habitation de Québec, built as a permanent fur trading outpost, was where he would forge a trading, and ultimately a military alliance, with the Algonquin and Huron nations.[79] First Nations traded their furs for many French goods such as metal objects, guns, alcohol, and clothing.
|
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|
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Coureurs des bois, voyageurs and Catholic missionaries used river canoes to explore the interior of the North American continent.[80] They established fur trading forts on the Great Lakes (Étienne Brûlé 1615), Hudson Bay (Radisson and Groseilliers 1659–60), Ohio River and Mississippi River (La Salle 1682), as well as the Saskatchewan River and Missouri River (de la Verendrye 1734–1738).[81]
|
71 |
+
|
72 |
+
After 1627, King Louis XIII of France allowed the Company of New France to introduced the seigneurial system and forbade settlement in New France by anyone other than Roman Catholics.[82]
|
73 |
+
|
74 |
+
In 1629 there was the surrender of Quebec, without battle, to English privateers led by David Kirke during the Anglo-French War. However, Samuel de Champlain argued that the English seizing of the lands was illegal as the war had already ended; he worked to have the lands returned to France. As part of the ongoing negotiations of their exit from the Anglo-French War, in 1632 the English king Charles agreed to return the lands in exchange for Louis XIII paying his wife's dowry. These terms were signed into law with the Treaty of Saint-Germain-en-Laye. The lands in Quebec and Acadia were returned to the French Company of One Hundred Associates.
|
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+
|
76 |
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New France became a Royal Province in 1663 under King Louis XIV of France with a Sovereign Council that included intendant Jean Talon.[83] The population grew slowly under French rule,[84] thus remained relatively low as growth was largely achieved through natural births, rather than by immigration.[85] To encourage population growth and to redress the severe imbalance between single men and women, King Louis XIV sponsored the passage of approximately 800 young French women (known as les filles du roi) to the colony. Most of the French were farmers ("Canadiens" or "Habitants"), and the rate of population growth among the settlers themselves was very high.[86]
|
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|
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Authorities in New France became more aggressive in their efforts to expel British traders and colonists from the Ohio Valley. They began construction of a series of fortifications to protect the area.[87] In 1754, George Washington launched a surprise attack on a group of Canadian soldiers sleeping in the early morning hours. It came at a time when no declaration of war had been issued by either country. This frontier aggression known as the Jumonville affair set the stage for the French and Indian War (a US designation; in Canada it is usually referred to as the Seven Years' War, although French Canadians often call it La guerre de la Conquête ["The War of Conquest"][88][89]) in North America. By 1756, France and Britain were battling the Seven Years' War worldwide. In 1758, the British mounted an attack on New France by sea and took the French fort at Louisbourg.
|
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On September 13, 1759, the British forces of General James Wolfe defeated those of French General Louis-Joseph de Montcalm on the Plains of Abraham outside Quebec City. With the exception of the small islands of Saint Pierre and Miquelon, located off the coast of Newfoundland, France ceded its North American possessions to Great Britain through the Treaty of Paris (1763) in favour of gaining the island of Guadeloupe for its then-lucrative sugar cane industry.[90] The British Royal Proclamation of 1763 renamed Canada (part of New France) as the Province of Quebec.
|
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|
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With unrest growing in the colonies to the south, which would one day grow into the American Revolution, the British were worried that the French-speaking Canadians might also support the growing rebellion. At that time, French-speaking Canadians formed the vast majority of the population of the province of Quebec (more than 99%) and British immigration was not going well. To secure the allegiance of the approximately 90,000 French-speaking Canadians to the British crown, first Governor James Murray and later Governor Guy Carleton promoted the need for change. There was also a need to compromise between the conflicting demands of the French-speaking Canadian subjects and those of newly arrived British subjects. These efforts by the colonial governors eventually resulted in enactment of the Quebec Act[91] of 1774.
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The Quebec Act provided the people of Quebec their first Charter of Rights and paved the way to later official recognition of the French language and French culture. The act also allowed the French speakers, known as Canadiens, to maintain French civil law and sanctioned freedom of religion, allowing the Roman Catholic Church to remain, one of the first cases in history of state-sanctioned freedom of religious practice.[92]
|
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+
|
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Although the Quebec Act was unrelated to the events in Boston of 1773, and was not regarded as one of the Coercive Acts, the timing of its passage led British colonists to the south to believe that it was part of the program to punish them. The Quebec Act offended a variety of interest groups in the British colonies. Land speculators and settlers objected to the transfer of western lands previously claimed by the colonies to a non-representative government. Many feared the establishment of Catholicism in Quebec, and that the French Canadians were being courted to help oppress British Americans.[93]
|
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|
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On June 27, 1775, General George Washington and his Continental Army invaded Canada in an attempt to conquer Quebec. British reinforcements came up the St. Lawrence in May 1776, and the Battle of Trois-Rivières turned into a disaster for the Americans. The army withdrew to Ticonderoga.[94] Although some help was given to the Americans by the locals, Governor Carleton punished American sympathizers, and public support of the American cause came to an end. In 1778, Frederick Haldimand took over for Guy Carleton as governor of Quebec.
|
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+
|
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The arrival of 10,000 Loyalists at Quebec in 1784 destroyed the political balance that Haldimand (and Carleton before him) had worked so hard to achieve. The swelling numbers of English encouraged them to make greater demands for recognition with the colonial government.[95] To restore stability to his largest remaining North American colony, King George III sent Carleton back to Quebec to remedy the situation.[96]
|
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|
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+
In ten years, Quebec had undergone a dramatic change. What worked for Carleton in 1774 was not likely to succeed in 1784. Specifically, there was no possibility of restoring the previous political balance – there were simply too many English people unwilling to reach a compromise with the 145,000 Canadiens or their colonial governor. The situation called for a more creative approach to problem solving.[96]
|
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+
|
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+
Loyalists soon petitioned the government to be allowed to use the British legal system they were
|
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+
used to in the American colonies. The creation of Upper and Lower Canada in 1791 allowed most Loyalists to live under British laws and institutions, while the French-speaking population of Lower Canada could maintain their familiar French civil law and the Catholic religion.[97] Therefore, Governor Haldimand (at the suggestion of Carleton) drew Loyalists away from Quebec City and Montreal by offering free land on the northern shore of Lake Ontario to anyone willing to swear allegiance to George III. The Loyalists were thus given land grants of 200 acres (81 ha) per person. Basically, this approach was designed with the intent of keeping French and English as far apart as possible. Therefore, after the separation of the Province of Quebec, Lower Canada and Upper Canada were formed, each with its own government.[96]
|
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In 1837, residents of Lower Canada – led by Louis-Joseph Papineau and Robert Nelson – formed an armed resistance group to seek an end to the unilateral control of the British governors.[98] They made a Declaration of Rights with equality for all citizens without discrimination and a Declaration of Independence of Lower Canada in 1838.[99] Their actions resulted in rebellions in both Lower and Upper Canada. An unprepared British Army had to raise militia force; the rebel forces scored a victory in Saint-Denis but were soon defeated.
|
98 |
+
|
99 |
+
After the rebellions, Lord Durham was asked to undertake a study and prepare a report on the matter and to offer a solution for the British Parliament to assess.[100] Following Durham's report,[100] the British government merged the two colonial provinces into one Province of Canada in 1840 with the Act of Union.[101] The two colonies remained distinct in administration, election, and law.
|
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+
|
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+
In 1848, Baldwin and LaFontaine, allies and leaders of the Reformist party, were asked by Lord Elgin to form an administration together under the new policy of responsible government. The French language subsequently regained legal status in the Legislature.[101]
|
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+
|
103 |
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In the 1860s, the delegates from the colonies of British North America (Canada, New Brunswick, Nova Scotia, Prince Edward Island, and Newfoundland) met in a series of conferences to discuss self-governing status for a new confederation. The first Charlottetown Conference took place in Charlottetown, Prince Edward Island, followed by the Quebec Conference in Quebec City which led to a delegation going to London, England, to put forth a proposal for a national union.[102]
|
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|
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As a result of those deliberations, in 1867 the Parliament of the United Kingdom passed the British North America Acts, providing for the Confederation of most of these provinces. The former Province of Canada was divided into its two previous parts as the provinces of Ontario (Upper Canada) and Quebec (Lower Canada). New Brunswick and Nova Scotia joined Ontario and Quebec in the new Dominion of Canada. The other provinces then joined the Confederation, one after the other: Manitoba and the Northwest Territories in 1870, British Columbia in 1871, Prince Edward Island in 1873, Yukon in 1898, Alberta and Saskatchewan in 1905, Newfoundland in 1949 and finally Nunavut in 1999.[103]
|
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|
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When Great Britain declared war on August 4, 1914, Canada was automatically involved as a dominion. About 6,000 volunteers from Quebec participated on the European front. Although reaction to conscription was favourable in English Canada the idea was deeply unpopular in Quebec. The Conscription Crisis of 1917 did much to highlight the divisions between French and English-speaking Canadians in Canada.
|
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|
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During World War II, the participation of Quebec was more important but led to the Conscription Crisis of 1944 and opposition. Many Quebecers fought against the axis powers between 1939 to 1945 with the involvement of many francophone regiments such as Les Fusiliers Mont-Royal, le Régiment de la Chaudière and many more.
|
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The conservative government of Maurice Duplessis and his Union Nationale dominated Quebec politics from 1944 to 1959 with the support of the Catholic Church.[104] Pierre Trudeau and other liberals formed an intellectual opposition to Duplessis's regime, setting the groundwork for the Quiet Revolution under Jean Lesage's Liberals. The Quiet Revolution was a period of dramatic social and political change[105] that saw the decline of Anglo supremacy in the Quebec economy, the decline of the Roman Catholic Church's influence,[14] the formation of hydroelectric companies under Hydro-Québec[105] and the emergence of a pro-sovereignty movement under former Liberal minister René Lévesque.
|
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Beginning in 1963, a paramilitary group that became known as the Front de libération du Québec (FLQ) launched a decade-long programme of propaganda and terrorism that included bombings, robberies and attacks[106] directed primarily at English institutions, resulting in at least five deaths. In 1970, their activities culminated in events referred to as the October Crisis when James Cross, the British trade commissioner to Canada, was kidnapped along with Pierre Laporte, a provincial minister and Vice-Premier.[107] Laporte was strangled with his own rosary beads a few days later. In their published Manifesto, the militants stated: "In the coming year Bourassa will have to face reality; 100,000 revolutionary workers, armed and organized." At the request of Premier Robert Bourassa, Prime Minister Pierre Trudeau invoked the War Measures Act.
|
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|
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In 1977, the newly elected Parti Québécois government of René Lévesque introduced the Charter of the French Language. Often known as Bill 101, it defined French as the only official language of Quebec in areas of provincial jurisdiction.[108]
|
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+
|
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Lévesque and his party had run in the 1970 and 1973 Quebec elections under a platform of separating Quebec from the rest of Canada. The party failed to win control of Quebec's National Assembly both times – though its share of the vote increased from 23 percent to 30 percent – and Lévesque was defeated both times in the riding he contested.[109] In the 1976 election campaign, he softened his message by promising a referendum (plebiscite) on sovereignty-association rather than outright separation, by which Quebec would have independence in most government functions but share some other ones, such as a common currency, with Canada. On November 15, 1976, Lévesque and the Parti Québécois won control of the provincial government for the first time. The question of sovereignty-association was placed before the voters in the 1980 Quebec referendum. During the campaign, Pierre Trudeau promised that a vote for the "no" side was a vote for reforming Canada. Trudeau advocated the patriation of Canada's Constitution from the United Kingdom. The existing constitutional document, the British North America Act, could only be amended by the United Kingdom Parliament upon a request by the Canadian parliament.
|
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|
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Sixty percent of the Quebec electorate voted against the proposition for sovereignty-association.[110] Polls showed that the overwhelming majority of English and immigrant Quebecers voted against, and that French Quebecers were almost equally divided, with older voters less in favour and younger voters more in favour. After his loss in the referendum, Lévesque went back to Ottawa to start negotiating a new constitution with Trudeau, his minister of Justice Jean Chrétien and the nine other provincial premiers. Lévesque insisted Quebec be able to veto any future constitutional amendments. The negotiations quickly reached a stand-still. Quebec is the only province not to have assented to the patriation of the Canadian constitution in 1982.[111]
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|
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+
In subsequent years, two attempts were made to gain Quebec's approval of the constitution. The first was the Meech Lake Accord of 1987, which was finally abandoned in 1990 when the province of Manitoba did not pass it within the established deadline. (Newfoundland premier Clyde Wells had expressed his opposition to the accord, but, with the failure in Manitoba, the vote for or against Meech never took place in his province.) This led to the formation of the sovereigntist Bloc Québécois party in Ottawa under the leadership of Lucien Bouchard,[112] who had resigned from the federal cabinet. The second attempt, the Charlottetown Accord of 1992, also failed to gain traction. This result caused a split in the Quebec Liberal Party that led to the formation of the new Action démocratique (Democratic Action) party led by Mario Dumont and Jean Allaire.
|
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+
|
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+
On October 30, 1995, with the Parti Québécois back in power since 1994, a second referendum on sovereignty took place. This time, it was rejected by a slim majority (50.6 percent NO to 49.4 percent YES).[113]
|
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|
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Given the province's heritage and the preponderance of French (unique among the Canadian provinces), there has been debate in Canada regarding the unique status (statut particulier) of Quebec and its people, wholly or partially. Prior attempts to amend the Canadian constitution to acknowledge Quebec as a "distinct society" – referring to the province's uniqueness within Canada regarding law, language, and culture – have been unsuccessful; however, the federal government under Prime Minister Jean Chrétien would later endorse recognition of Quebec as a distinct society.[114]
|
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On October 30, 2003, the National Assembly of Quebec voted unanimously to affirm "that the people of Québec form a nation".[115] On November 27, 2006, the House of Commons passed a symbolic motion moved by Prime Minister Stephen Harper declaring "that this House recognize that the Québécois form a nation within a united Canada."[116][117][118] However, there is considerable debate and uncertainty over what this means.[119][120]
|
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+
|
129 |
+
The Lieutenant Governor represents the Queen of Canada and acts as the province's head of state.[121][122] The head of government is the premier (called premier ministre in French) who leads the largest party in the unicameral National Assembly, or Assemblée Nationale, from which the Executive Council of Quebec is appointed.
|
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+
|
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+
Until 1968, the Quebec legislature was bicameral,[123] consisting of the Legislative Council and the Legislative Assembly. In that year, the Legislative Council was abolished and the Legislative Assembly was renamed the National Assembly. Quebec was the last province to abolish its legislative council.
|
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|
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The government of Quebec awards an order of merit called the National Order of Quebec. It is inspired in part by the French Legion of Honour. It is conferred upon men and women born or living in Quebec (but non-Quebecers can be inducted as well) for outstanding achievements.[124]
|
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|
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The government of Quebec takes the majority of its revenue through a progressive income tax, a 9.975% sales tax[125] and various other taxes (such as carbon, corporate and capital gains taxes), equalization payments from the federal government, transfer payments from other provinces and direct payments.[126] By some measures Quebec is the highest taxed province;[127] a 2012 study indicated that "Quebec companies pay 26 per cent more in taxes than the Canadian average".[128] A 2014 report by the Fraser Institute indicated that "Relative to its size, Quebec is the most indebted province in Canada by a wide margin".[129]
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|
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+
Quebec has subdivisions at the regional, supralocal and local levels. Excluding administrative units reserved for Aboriginal lands, the primary types of subdivision are:
|
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+
|
139 |
+
At the regional level:
|
140 |
+
|
141 |
+
At the supralocal level:
|
142 |
+
|
143 |
+
At the local level:
|
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+
|
145 |
+
In the 2016 census, Quebec had a population of 8,164,361 living in 3,531,663 of its 3,858,943 total dwellings, a 3.3% change from its 2011 population of 7,903,001. With a land area of 1,356,625.27 km2 (523,795.95 sq mi), it had a population density of 6.0/km2 (15.6/sq mi) in 2016.[1] In 2013, Statistics Canada estimated the province's population to be 8,155,334.[132]
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At 1.69 children per woman, Quebec's 2011 fertility rate is above the Canada-wide rate of 1.61,[133] and is higher than it was at the turn of the 21st century. However, it is still below the replacement fertility rate of 2.1. This contrasts with its fertility rates before 1960, which were among the highest of any industrialized society. Although Quebec is home to only 24% of the population of Canada, the number of international adoptions in Quebec is the highest of all provinces of Canada. In 2001, 42% of international adoptions in Canada were carried out in Quebec. By 2012, the population of Quebec reached 8 million, and it is projected to reach 9.2 million in 2056.[134] Life expectancy in Quebec reached a new high in 2011, with an expectancy of 78.6 years for men and 83.2 years for women; this ranked as the third-longest life expectancy among Canadian provinces, behind those of British Columbia and Ontario.[133]
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All the tables in the following section have been reduced from their original size, for full tables see main article Demographics of Quebec.
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+
|
151 |
+
Origins in this table are self-reported and respondents were allowed to give more than one answer.
|
152 |
+
|
153 |
+
Percentages are calculated as a proportion of the total number of respondents (7,435,905) and may total more than 100 percent due to dual responses. Only groups with 1.5 percent or more of respondents are shown.[135]
|
154 |
+
|
155 |
+
The 2006 census counted a total aboriginal population of 108,425 (1.5 percent) including 65,085 North American Indians (0.9 percent), 27,985 Métis (0.4 percent), and 10,950 Inuit (0.15 percent). There is a significant undercount, as many of the largest Indian bands regularly refuse to participate in Canadian censuses for political reasons regarding the question of aboriginal sovereignty. In particular, the largest Mohawk Iroquois reserves (Kahnawake, Akwesasne and Kanesatake) were not counted.
|
156 |
+
|
157 |
+
Nearly 9% of the population of Quebec belongs to a visible minority group. This is a lower percentage than that of British Columbia, Ontario, Alberta, and Manitoba but higher than that of the other five provinces. Most visible minorities in Quebec live in or near Montreal.
|
158 |
+
|
159 |
+
Percentages are calculated as a proportion of the total number of respondents (7,435,905).Only groups with more than 0.5 percent of respondents are shown.[136]
|
160 |
+
|
161 |
+
Religion in Quebec (2011 National Household Survey):[137]
|
162 |
+
|
163 |
+
Quebec is unique among the provinces in its overwhelmingly Roman Catholic population, though recently with a low church attendance. This is a legacy of colonial times when only Roman Catholics were permitted to settle in New France. The 2001 census showed the population to be 90.3 percent Christian (in contrast to 77 percent for the whole country) with 83.4 percent Catholic (including 83.2 percent Roman Catholic); 4.7 percent Protestant Christian (including 1.2 percent Anglican, 0.7 percent United Church; and 0.5 percent Baptist); 1.4 percent Orthodox Christian (including 0.7 percent Greek Orthodox); and 0.8 percent other Christian; as well as 1.5 percent Muslim; 1.3 percent Jewish; 0.6 percent Buddhist; 0.3 percent Hindu; and 0.1 percent Sikh. An additional 5.8 percent of the population said they had no religious affiliation (including 5.6 percent who stated that they had no religion at all).Percentages are calculated as a proportion of the total number of respondents (7,125,580)[138]
|
164 |
+
|
165 |
+
The official language of Quebec is French. Quebec is the only Canadian province whose population is mainly Francophone; 6,102,210 people (78.1 percent of the population) recorded it as their sole native language in the 2011 Census, and 6,249,085 (80.0%) recorded that they spoke it most often at home.[139] Knowledge of French is widespread even among those who do not speak it natively; in 2011, about 94.4 percent of the total population reported being able to speak French, alone or in combination with other languages, while 47.3% reported being able to speak English.[139]
|
166 |
+
|
167 |
+
In 2011, 599,230 people (7.7 percent of the population) people in Quebec declared English to be their mother tongue, and 767,415 (9.8 percent) used it most often as their home language[139] The English-speaking community or Anglophones are entitled to services in English in the areas of justice, health, and education;[140] services in English are offered in municipalities in which more than half the residents have English as their mother tongue. Allophones, people whose mother tongue is neither French nor English, made up 12.3 percent (961,700) of the population, according to the 2011 census, though a smaller figure – 554,400 (7.1 percent) – actually used these languages most often in the home.[139]
|
168 |
+
|
169 |
+
A considerable number of Quebec residents consider themselves to be bilingual in French and English. In Quebec, about 42.6 percent of the population (3,328,725 people) report knowing both languages; this is the highest proportion of bilinguals of any Canadian province.[139] One specific area in the Bilingual Belt called the West Island of Montreal, represented by the federal electoral district of Lac-Saint-Louis, is the most bilingual area in the province: 72.8% of its residents claim to know English and French according to the most recent census.[141] In contrast, in the rest of Canada, in 2006 only about 10.2 percent (2,430,990) of the population had a knowledge of both of the country's official languages. Altogether, 17.5% of Canadians are bilingual in French and English.[139]
|
170 |
+
|
171 |
+
In 2011, the most common mother tongue languages in the province were as follows:
|
172 |
+
(Figures shown are for single-language responses only.)
|
173 |
+
|
174 |
+
Following were Creoles (0.8%), Chinese (0.6%), Greek (0.5%), Portuguese (0.5%), Romanian (0.4%), Vietnamese (0.3%), and Russian (0.3%). In addition, 152,820 (2.0%) reported having more than one native language.[139]
|
175 |
+
|
176 |
+
English is not designated an official language by Quebec law.[140] However, both English and French are required by the Constitution Act, 1867, for the enactment of laws and regulations, and any person may use English or French in the National Assembly and the courts. The books and records of the National Assembly must also be kept in both languages.[142][143] Until 1969, Quebec was the only officially bilingual province in Canada and most public institutions functioned in both languages. English was also used in the legislature, government commissions and courts.
|
177 |
+
|
178 |
+
Since the 1970s, languages other than French on commercial signs have been permitted only if French is given marked prominence. This law has been the subject of periodic controversy since its inception. The written forms of French place-names in Canada retain their diacritics such as accent marks over vowels in English text. Legitimate exceptions are Montreal and Quebec. However, the accented forms are increasingly evident in some publications. The Canadian Style states that Montréal and Québec (the city) must retain their accents in English federal documents.
|
179 |
+
|
180 |
+
Montreal
|
181 |
+
Quebec City
|
182 |
+
Gatineau
|
183 |
+
|
184 |
+
Quebec has an advanced, market-based, and open economy. In 2009, its gross domestic product (GDP) of US$32,408 per capita at purchasing power parity puts the province at par with Japan, Italy and Spain, but remains lower than the Canadian average of US$37,830 per capita.[145][verification needed] The economy of Quebec is ranked the 37th largest economy in the world just behind Greece and 28th for the gross domestic product (GDP) per capita.[146][147]
|
185 |
+
|
186 |
+
The economy of Quebec represents 20.36% of the total GDP of Canada. Like most industrialized countries, the economy of Quebec is based mainly on the services sector. Quebec's economy has traditionally been fuelled by abundant natural resources, a well-developed infrastructure, and average productivity. The provincial GDP in 2010 was C$319,348 billion,[148] which makes Quebec the second largest economy in Canada.
|
187 |
+
|
188 |
+
The provincial debt-to-GDP ratio peaked at 50.7% in fiscal year 2012–2013, and is projected to decline to 33.8% in 2023–2024.[149] The credit rating of Quebec is currently Aa2 according to the Moody's agency.[150] In June 2017 S&P rated Quebec as an AA- credit risk, surpassing Ontario for the first time.[151]
|
189 |
+
|
190 |
+
Quebec's economy has undergone tremendous changes over the last decade.[152] Firmly grounded in the knowledge economy, Quebec has one of the highest growth rate of gross domestic product (GDP) in Canada. The knowledge sector represents about 30.9% of Quebec's GDP.[153] Quebec is experiencing faster growth of its R&D spending than other Canadian provinces.[154] Quebec's spending in R&D in 2011 was equal to 2.63% of GDP, above the European Union average of 1.84% and will have to reaches the target of devoting 3% of GDP to research and development activities in 2013 according to the Lisbon Strategy.[155] The percentage spent on research and technology (R&D) is the highest in Canada and higher than the averages for the Organisation for Economic Co-operation and Development and the G7 countries.[156] Approximately 1.1 million Quebecers work in the field of science and technology.[157]
|
191 |
+
|
192 |
+
Quebec is also a major player in several leading-edge industries including aerospace, information technologies and software and multimedia. Approximately 60% of the production of the Canadian aerospace industry are from Quebec, where sales totalled C$12.4 billion in 2009.[161] Quebec is one of North America's leading high-tech player. This vast sector encompassing approximately 7,300 businesses and employ more than 145,000 people.[162] Pauline Marois has recently unveiled a two billion dollar budget for the period between 2013 to 2017 to create about 115,000 new jobs in knowledge and innovation sectors. The government promises to provide about 3% of Quebec's GDP in research and development (R&D).[163]
|
193 |
+
|
194 |
+
About 180 000 Quebeckers work in different field of information technology.[164] Approximately 52% of Canadian companies in these sectors are based in Quebec, mainly in Montreal and Quebec City. There are currently approximately 115 telecommunications companies established in the province, such as Motorola and Ericsson. About 60 000 people currently working in computer software development. Approximately 12 900 people working in over 110 companies such as IBM, CMC, and Matrox. The multimedia sector is also dominated by the province of Quebec. Several companies, such as Ubisoft settled in Quebec since the late 1990s.[165]
|
195 |
+
|
196 |
+
The mining industry accounted for 6.3% of Quebec's GDP.[166] It employs about 50,000 people[167] in 158 companies.[167]
|
197 |
+
|
198 |
+
The pulp and paper industries generate annual shipments valued at more than $14 billion.[168] The forest products industry ranks second in exports, with shipments valued at almost $11 billion. It is also the main, and in some circumstances only, source of manufacturing activity in more than 250 municipalities in the province. The forest industry has slowed in recent years because of the softwood lumber dispute.[169] This industry employs 68,000 people in several regions of Quebec.[170] This industry accounted for 3.1% of Quebec's GDP.[171]
|
199 |
+
|
200 |
+
Agri-food industry plays an important role in the economy of Quebec, with meat and Dairy products being the two main sectors. It accounts for 8% of the Quebec's GDP and generate $19.2 billion. This industry generated 487,000 jobs in agriculture, fisheries, manufacturing of food, beverages and tobacco and food distribution.[172]
|
201 |
+
|
202 |
+
The abundance of natural resources gives Quebec an advantageous position on the world market. Quebec stands out particularly in the mining sector, ranking among the top ten areas to do business in mining.[173] It also stands for the exploitation of its forest resources.
|
203 |
+
|
204 |
+
Quebec is remarkable for the natural resources of its vast territory. It has about 30 mines, 158 exploration companies and fifteen primary processing industries. Many metallic minerals are exploited, the principals are gold, iron, copper and zinc. Many other substances are extracted including titanium, asbestos, silver, magnesium, nickel and many other metals and industrial minerals.[174] However, only 40% of the mineral potential of Quebec is currently known. In 2003, the value of mineral exploitation reached Quebec 3.7 billion Canadian dollars.[175] Moreover, as a major centre of exploration for diamonds,[176] Quebec has seen, since 2002, an increase in its mineral explorations, particularly in the Northwest as well as in the Otish Mountains and the Torngat Mountains.
|
205 |
+
|
206 |
+
The vast majority (90.5%) of Quebec's forests are publicly owned. Forests cover more than half of Quebec's territory, for a total area of nearly 761,100 square kilometres (293,900 sq mi).[177] The Quebec forest area covers seven degrees of latitude.
|
207 |
+
|
208 |
+
More than a million lakes and rivers cover Quebec, occupying 21% of the total area of its territory. The aquatic environment is composed of 12.1% of fresh water and 9.2% of saltwater (percentage of total QC area).[178]
|
209 |
+
|
210 |
+
The government of Quebec has launched the Stratégie québécoise de la recherche et de l'innovation (SQRI) in 2007 which aims to promote development through research, science and technology. The government hopes to create a strong culture of innovation in Quebec for the next decades and to create a sustainable economy.[179] The spending on research and development reached some 7.824 billion dollars in 2007, roughly the equivalent of 2.63% of Quebec's GDP.[179] Quebec is ranked, as of March 2011, 13th in the world in terms of investment in research and development.[180] The research and development expenditures will be more than 3% of the province's GDP in 2013. The R&D expenditure in Quebec is higher than the average G7 and OECD countries.[157] Science and technology are key factors in the economic position of Quebec. More than one million people in Quebec are employed in the science and technology sector.[157]
|
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+
|
212 |
+
Quebec is considered as one of world leaders in fundamental scientific research, having produced ten Nobel laureates in either physics, chemistry, or medicine.[182] It is also considered as one of the world leaders in sectors such as aerospace, information technology, biotechnology and pharmaceuticals, and therefore plays a significant role in the world's scientific and technological communities.[183] Quebec is also active in the development of its energy industries, including renewable energy such as hydropower and wind power. Quebec has had over 9,469 scientific publications in the sector of medicine, biomedical research and engineering since the year 2000.[184] Overall, the province of Quebec count about 125 scientific publications per 100,000 inhabitants in 2009.[185] The contribution of Quebec in science and technology represent approximately 1% of the researches worldwide since the 1980s to 2009.[186] Between 1991 to 2000, Quebec produced more scientific papers per 100,000 inhabitants than the United States and Germany.[187]
|
213 |
+
|
214 |
+
The Canadian Space Agency was established in Quebec due to its major role in this research field. A total of three Quebecers have been in space since the creation of the CSA: Marc Garneau, Julie Payette and Guy Laliberté. Quebec has also contributed to the creation of some Canadian artificial satellites including SCISAT-1, ISIS, Radarsat-1 and Radarsat-2.[188][189][190]
|
215 |
+
|
216 |
+
The province is one of the world leaders in the field of space science and contributed to important discoveries in this field.[191] One of the most recent is the discovery of the complex extrasolar planets system HR 8799. HR 8799 is the first direct observation of an exoplanet in history.[192][193] Olivier Daigle and Claude Carignan, astrophysicists from Université de Montréal have invented an astronomical camera approximately 500 times more powerful than those currently on the market.[194] It is therefore considered as the most sensitive camera in the world.[195][196][197] The Mont Mégantic Observatory was recently equipped with this camera.[198]
|
217 |
+
|
218 |
+
Quebec ranks among the world leaders in the field of life science.[199] William Osler, Wilder Penfield, Donald Hebb, Brenda Milner, and others made significant discoveries in medicine, neuroscience and psychology while working at McGill University in Montreal. Quebec has more than 450 biotechnology and pharmaceutical companies which together employ more than 25,000 people and 10,000 highly qualified researchers.[199] Montreal is ranked 4th in North America for the number of jobs in the pharmaceutical sector.[199][200]
|
219 |
+
|
220 |
+
Development and security of land transportation in Canada are provided by the ministère des Transports du Québec.[201] Other organizations, such as the Canadian Coast Guard and Nav Canada, provide the same service for the sea and air transportation. The Commission des transports du Québec works with the freight carriers and the public transport.
|
221 |
+
|
222 |
+
The réseau routier québécois (Quebec road network) is managed by the Société de l'assurance automobile du Québec (SAAQ) (Quebec Automobile Insurance Corporation) and consists of about 185,000 kilometres (115,000 mi) of highways and national, regional, local, collector and forest roads. In addition, Quebec has almost 12,000 bridges, tunnels, retaining walls, culverts and other structures[202] such as the Quebec Bridge, the Laviolette Bridge and the Louis-Hippolyte Lafontaine Bridge–Tunnel.
|
223 |
+
|
224 |
+
In the waters of the St. Lawrence there are eight deep-water ports for the transhipment of goods. In 2003, 3886 cargo and 9.7 million tonnes of goods transited the Quebec portion of the St. Lawrence Seaway.[203]
|
225 |
+
|
226 |
+
Concerning rail transport, Quebec has 6,678 kilometres (4,150 mi) of railways[204] integrated in the large North American network. Although primarily intended for the transport of goods through companies such as the Canadian National (CN) and the Canadian Pacific (CP), the Quebec railway network is also used by inter-city passengers via Via Rail Canada and Amtrak. In April 2012, plans were unveiled for the construction of an 800 km (497 mi) railway running north from Sept-Îles, to support mining and other resource extraction in the Labrador Trough.[205]
|
227 |
+
|
228 |
+
The upper air network includes 43 airports that offer scheduled services on a daily basis.[203] In addition, the Government of Quebec owns airports and heliports to increase the accessibility of local services to communities in the Basse-Côte-Nord and northern regions.[206]
|
229 |
+
|
230 |
+
Various other transport networks crisscross the province of Quebec, including hiking trails, snowmobile trails and bike paths; the Green Road being the largest with nearly 4,000 kilometres (2,500 mi) in length.[207]
|
231 |
+
|
232 |
+
Quebec has been described as a potential clean energy superpower.[208][209] The energy balance of Quebec has undergone a large shift over the past 30 years. In 2008, electricity ranked as the main form of energy used in Quebec (41.6%), followed by oil (38.2%) and natural gas (10.7%).[210]
|
233 |
+
|
234 |
+
Quebec is the fourth largest producer of hydroelectricity in the world after China, Brazil and the United States and relies almost exclusively (96% in 2008) on this source of renewable energy for its electricity needs.[211]
|
235 |
+
|
236 |
+
Quebec is at the centre of French-speaking culture in North America. Its culture is a symbol of a distinct perspective. Quebec nationalism has been one expression of this perspective. Quebec's culture blends its historic roots with its aboriginal heritage and the contributions of recent immigrants, as well as receiving a strong influence from English-speaking North America.
|
237 |
+
|
238 |
+
Montreal's cabarets rose to the forefront of the city's cultural life during the Prohibition era of Canada and the United States in the 1920s. The cabarets radically transformed the artistic scene, greatly influencing the live entertainment industry of Quebec.[212] The Quartier Latin (English: Latin Quarter) of Montreal, and Vieux-Québec (English: Old Quebec) in Quebec City, are two hubs of activity for today's artists. Life in the cafés and "terrasses" (outdoor restaurant terraces) reveals a Latin influence in Quebec's culture, with the théâtre Saint-Denis in Montréal and the Capitole de Québec theatre in Quebec City being among the principal attractions.
|
239 |
+
|
240 |
+
A number of governmental and non-government organizations support cultural activity in Quebec. The Conseil des arts et des lettres du Québec (CALQ) is an initiative of the Ministry of Culture and Communications (Quebec). It supports creation, innovation, production, and international exhibits for all cultural fields of Quebec. The Société de développement des entreprises culturelles (SODEC) works to promote and fund individuals working in the cultural industry. The Prix du Québec is an award given by the government to confer the highest distinction and honour to individuals demonstrating exceptional achievement in their respective cultural field.
|
241 |
+
|
242 |
+
On February 8, 2007, Quebec Premier Jean Charest announced the setting up of a Commission tasked with consulting Quebec Society on the matter of arrangements regarding cultural diversity. The Premier's press release[213] reasserted the three fundamental values of Quebec society:
|
243 |
+
|
244 |
+
Equality between men and women, primacy of the French language, and separation of church and state constitute the fundamental values. They are not subject to any arrangement. They cannot be subordinated by any other principle.[213]
|
245 |
+
|
246 |
+
Furthermore, Quebec is a free and democratic society that abides by the rule of law.[214] Quebec society bases its cohesion and specificity on a set of statements, a few notable examples of which include:
|
247 |
+
|
248 |
+
Traditional music is imbued with many dances, such as the jig, the quadrille, the reel and line dancing, which developed in the festivities since the early days of colonization. Various instruments are more popular in Quebec's culture: harmonica (music-of-mouth or lip-destruction), fiddle, spoons, jaw harp and accordion. The podorythmie is a characteristic of traditional Quebec music and means giving the rhythm with the feet.[218] Quebec traditional music is currently provided by various contemporary groups seen mostly during Christmas and New Year's Eve celebrations, Quebec National Holiday and many local festivals.
|
249 |
+
|
250 |
+
Being a modern cosmopolitan society, today, all types of music can be found in Quebec. From folk music to hip-hop, music has always played an important role in Quebercers culture. From La Bolduc in the 1920s–1930s to the contemporary artists, the music in Quebec has announced multiple songwriters and performers, pop singers and crooners, music groups and many more. Quebec's most popular artists of the last century include the singers Félix Leclerc (1950s), Gilles Vigneault (1960s–present), Kate and Anna McGarrigle (1970s–present) and Céline Dion (1980s–present).[219] The First Nations and the Inuit of Quebec also have their own traditional music.
|
251 |
+
|
252 |
+
From Quebec's musical repertoire, the song A La Claire Fontaine[220] was the anthem of the New France, Patriots and French Canadian, then replaced by O Canada. Currently, the song Gens du pays is by far preferred by many Quebecers to be the national anthem of Quebec. The Association québécoise de l'industrie du disque, du spectacle et de la vidéo (ADISQ) was created in 1978 to promote the music industry in Quebec.[221] The Orchestre symphonique de Québec and the Orchestre symphonique de Montréal are respectively associated with the Opéra de Québec and the Opéra de Montreal whose performances are presented at the Grand Théâtre de Québec and at Place des Arts. The Ballets Jazz de Montreal, the Grands Ballets and La La La Human Steps are three important professional troupes of contemporary dance.
|
253 |
+
|
254 |
+
The Cinémathèque québécoise has a mandate to promote the film and television heritage of Quebec. Similarly, the National Film Board of Canada (NFB), a federal Crown corporation, provides for the same mission in Canada.
|
255 |
+
In a similar way, the Association of Film and Television in Quebec (APFTQ) promotes independent production in film and television.[222] While the Association of Producers and Directors of Quebec (APDQ) represents the business of filmmaking and television, the Association of Community Radio Broadcasters of Quebec (ARCQ) (French acronym) represents the independent radio stations.[223] Several movie theatres across Quebec ensure the dissemination of Quebec cinema. With its cinematic installations, such as the Cité du cinéma and Mel's studios, the city of Montreal is home to the filming of various productions.[224] The State corporation Télé-Québec, the federal Crown corporation CBC, general and specialized private channels, networks, independent and community radio stations broadcast the various Quebec téléromans, the national and regional news, interactive and spoken programmations, etc.[225][226]
|
256 |
+
Les Rendez-vous du cinéma québécois is a festival surrounding the ceremony of the Jutra Awards Night that rewards work and personalities of Quebec cinema.[227] The Artis and the Gemini Awards gala recognize the personalities of television and radio industry in Quebec and French Canada. The Film Festival of the 3 Americas, Quebec City, the Festival of International Short Film, Saguenay, the World Film Festival and the Festival of New Cinema, Montreal, are other annual events surrounding the film industry in Quebec.
|
257 |
+
|
258 |
+
From New France, Quebec literature was first developed in the travel accounts of explorers such as Jacques Cartier, Jean de Brébeuf, the Baron de La Hontan and Nicolas Perrot, describing their relations with indigenous peoples. The Moulin à paroles traces the great texts that have shaped the history of Quebec since its foundation in 1534 until the era of modernity. The first to write the history of Quebec, since its discovery, was the historian François-Xavier Garneau. This author will be part of the current of patriotic literature (also known as the "poets of the country" and literary identity) that will arise after the Patriots Rebellion of 1837–1838.[228]
|
259 |
+
|
260 |
+
Various tales and stories are told through oral tradition, such as, among many more, the legends of the Bogeyman, the Chasse-galerie, the Black Horse of Trois-Pistoles, the Complainte de Cadieux, the Corriveau, the dancing devil of Saint-Ambroise, the Giant Beaupré, the monsters of the lakes Pohénégamook and Memphremagog, of Quebec Bridge (called the Devil's Bridge), the Rocher Percé and of Rose Latulipe, for example.[229]
|
261 |
+
|
262 |
+
Many Quebec poets and prominent authors marked their era and today remain anchored in the collective imagination, like, among others, Philippe Aubert de Gaspé, Octave Crémazie, Honoré Beaugrand, Émile Nelligan, Lionel Groulx, Gabrielle Roy, Hubert Aquin, Michel Tremblay, Marie Laberge, Fred Pellerin and Gaston Miron. The regional novel from Quebec is called Terroir novel and is a literary tradition[230] specific to the province. It includes such works as The Old Canadians, Maria Chapdelaine, Un homme et son péché, Le Survenant, etc. There are also many successful plays from this literary category, such as Les Belles-sœurs and Broue (Brew).
|
263 |
+
|
264 |
+
Among the theatre troupes are the Compagnie Jean-Duceppe, the Théâtre La Rubrique at the Pierrette-Gaudreault venue of the Institut of arts in Saguenay, the Théâtre Le Grenier, etc. In addition to the network of cultural centres in Quebec,[231] the venues include the Monument-National and the Rideau Vert (green curtain) Theatre in Montreal, the Trident Theatre in Quebec City, etc. The National Theatre School of Canada and the Conservatoire de musique et d'art dramatique du Québec form the future players.
|
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|
266 |
+
Popular French-language contemporary writers include Louis Caron, Suzanne Jacob, Yves Beauchemin, and Gilles Archambault. Mavis Gallant, born in Quebec, lived in Paris from the 1950s onward. Well-known English-language writers from Quebec include Leonard Cohen, Mordecai Richler, and Neil Bissoondath.
|
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+
|
268 |
+
Henri Julien, La Chasse-galerie
|
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+
|
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+
Émile Nelligan, a Quebec poet famous for his poem Winter evening
|
271 |
+
|
272 |
+
First influenced since the days of New France by Catholicism, with works from Frère Luc (Brother Luke) and more recently from Ozias Leduc and Guido Nincheri, art of Quebec has developed around the specific characteristics of its landscapes and cultural, historical, social and political representations.[232]
|
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+
|
274 |
+
Thus, the development of Quebec masterpieces in painting, printmaking and sculpture is marked by the contribution of artists such as Louis-Philippe Hébert, Cornelius Krieghoff, Alfred Laliberté, Marc-Aurèle Fortin, Marc-Aurèle de Foy Suzor-Coté, Jean Paul Lemieux, Clarence Gagnon, Adrien Dufresne, Alfred Pellan, Jean-Philippe Dallaire, Charles Daudelin, Arthur Villeneuve, Jean-Paul Riopelle, Paul-Émile Borduas and Marcelle Ferron.
|
275 |
+
|
276 |
+
The Fine arts of Quebec are displayed at the Quebec National Museum of Fine Arts, the Montreal Museum of Contemporary Art, the Montreal Museum of Fine Arts, the Quebec Salon des métiers d'art and in many art galleries. While many works decorate the public areas of Quebec, others are displayed in foreign countries such as the sculpture Embâcle (Jam) by Charles Daudelin on Québec Place in Paris and the statue Québec Libre! (free Quebec!) by Armand Vaillancourt in San Francisco.
|
277 |
+
The Montreal School of Fine Arts forms the painters, printmakers and sculptors of Quebec.
|
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+
|
279 |
+
Various buildings reflect the architectural heritage that characterizes Quebec, such as religious buildings, city halls, houses of large estates, and other locations throughout the province.
|
280 |
+
|
281 |
+
Several circus troupes were created in recent decades, the most important being without any doubt the Cirque du Soleil.[233] Among these troops are contemporary, travelling and on-horseback circuses, such as Les 7 Doigts de la Main, Cirque Éloize, Cavalia, Kosmogonia, Saka and Cirque Akya.[234] Presented outdoors under a tent or in venues similar to the Montreal Casino, the circuses attract large crowds both in Quebec and abroad. In the manner of touring companies of the Renaissance, the clowns, street performers, minstrels, or troubadours travel from city to city to play their comedies. Although they may appear randomly from time to time during the year, they are always visible in the cultural events such as the Winterlude in Gatineau, the Quebec Winter Carnival, the Gatineau Hot Air Balloon Festival, the Quebec City Summer Festival, the Just for Laughs Festival in Montreal and the Festival of New France in Quebec.
|
282 |
+
|
283 |
+
The National Circus School and the École de cirque de Québec were created to train future Contemporary circus artists. For its part, Tohu, la Cité des Arts du Cirque was founded in 2004 to disseminate the circus arts.[235]
|
284 |
+
|
285 |
+
The Cultural Heritage Fund is a program of the Quebec government[236] for the conservation and development of Quebec's heritage, together with various laws.[237] Several organizations ensure that same mission, both in the social and cultural traditions in the countryside and heritage buildings, including the Commission des biens culturels du Québec, the Quebec Heritage Fondation, the Conservation Centre of Quebec, the Centre for development of living heritage, the Quebec Council of living heri tage, the Quebec Association of heritage interpretation, etc.
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|
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Several sites, houses and historical works reflect the cultural heritage of Quebec, such as the Village Québécois d'Antan, the historical village of Val-Jalbert, the Fort Chambly, the national home of the Patriots, the Chicoutimi pulp mill (Pulperie de Chicoutimi), the Lachine Canal and the Victoria Bridge. Strongly influenced by the presence of the Catholic Church, the development of the religious history of Quebec is provided by organizations like the Council of the religious heritage of Quebec. Since 2007, the government promotes, with the various players in the field, the conclusion of agreements on the use of property belonging to episcopal factories and corporations to establish "partnerships in financing the restoration and renovation of religious buildings".[238]
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As of December 2011, there are 190 National Historic Sites of Canada in Quebec.[239] These sites were designated as being of national historic significance.[240]
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Various museums tell the cultural history of Quebec, like the Museum of Civilization, the Museum of French America, the McCord Museum or the Montreal Museum of Archaeology and History in Pointe-à-Callière, displaying artifacts, paintings and other remains from the past of Quebec. Many literary works reproduce the daily lives of the past, following the social and cultural traditions of Quebec television series reproducing the old days[241] such as the trilogy of Pierre Gauvreau (Le Temps d'une paix, Cormoran and Le Volcan tranquille), La Famille Plouffe, Les Belles Histoires des Pays-d'en-Haut, La Petite Patrie, Entre chien et loup, Les Filles de Caleb, Blanche, Au nom du père et du fils, Marguerite Volant, Nos Étés or Musée Éden, among others.
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The historical context of 'traditional' Quebec cuisine is from the fur trade period and many dishes have a high fat or lard content. From the early 17th century, French settlers populating North America were interested in a new cuisine to confront the climate and the needs arising from the work of colonization. Mindful of the same nutritional needs as settlers from Acadia, it has many similarities with Acadian cuisine. Quebec's cuisine has a strong French and Irish influence, although many aspects of Canadian aboriginal cuisine have also had a significant impact on Quebec cuisine. Quebec is most famous for its Tourtière, Pâté Chinois, Poutine, St. Catherine's taffy among others. The temps des sucres (sugar season) is one of the oldest of Quebec culinary traditions. During springtime, many Quebecers go to the cabane à sucre (sugar house) for a traditional meal. The Jewish community of Montreal has contributed Montreal-style bagels and smoked meat which is similar to pastrami.
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Quebec has produced beer since the beginning of colonization especially with the emergence of spruce beer. In 1668, Jean Talon founded a brewery in Quebec City, but it closed a decade later. Although many people tried to produce a beer between the 17th and 18th centuries, it is only since the 1980s that the industries had produced on a larger scale. Today there are nearly a hundred breweries and companies, including Unibroue, Molson Coors, Labatt and many others. Quebec also produces wine, ice wine and ice cider.
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Quebec has produced cheese for centuries. The first cheese-making school in North America was established in Saint-Denis-de-Kamouraska in 1893. It was at this moment that the monks of La Trappe of Oka began to produce the famous Oka cheese. Today there are over 300 different cheeses in Quebec.
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Sports in Quebec constitutes an essential dimension of Quebec culture. The practice of sports and outdoor activities in Quebec was influenced largely by its geography and climate. Ice hockey remains the national sport. This sport, which was played for the first time on March 3, 1875, at the Victoria Skating Rink in Montreal and promoted over the years by numerous achievements, including the centenary of the Montreal Canadiens, still raises passions.[242] Other major sports include Canadian football with the Montreal Alouettes, soccer with the Montreal Impact, the Grand Prix du Canada Formula 1 racing with drivers such as Gilles Villeneuve and Jacques Villeneuve, and professional baseball with the former Montreal Expos. During its history, Quebec has hosted several major sporting events; including the 1976 Summer Olympics, the Fencing World Championships in 1967, track cycling in 1974, and the Transat Québec-Saint-Malo race created for the first time in 1984.
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Québec athletes have performed well at the Winter Olympics over recent years. They won 12 of Canada's 29 medals at the most recent Winter Olympics in Pyeongchang (2018); they won 12 of the 27 Canadian medals in Sochi (2014); and 9 of the 26 Canadian medals in Vancouver (2010).[243]
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In 1939, the government of Quebec unilaterally ratified its coat of arms to reflect Quebec's political history: French rule (gold lily on blue background), British rule (lion on red background) and Canadian rule (maple leaves) and with Quebec's motto below "Je me souviens".[244] Je me souviens ("I remember") was first carved under the coat of arms of Quebec's Parliament Building façade in 1883. It is an official part of the coat of arms and has been the official licence plate motto since 1978, replacing "La belle province" (the beautiful province). The expression La belle province is still used mostly in tourism as a nickname for the province.[citation needed]
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The fleur-de-lis, the ancient symbol of the French monarchy, first arrived on the shores of the Gaspésie in 1534 with Jacques Cartier on his first voyage. When Samuel de Champlain founded Québec City in 1608, his ship hoisted the merchant flag of a white cross on a blue background. By 1758 at the Battle of Carillon, the Flag of Carillon would become the basis of Quebec's desire to have its own flag. By 1903, the parent of today's flag had taken shape, known as the "Fleurdelisé". The flag in its present form with its 4 white "fleur-de-lis" lilies on a blue background with a white cross replaced the Union Jack on Quebec's Parliament Building on January 21, 1948.
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In 1998 the Montreal Insectarium sponsored a poll to choose an official insect. The white admiral butterfly (Limenitis arthemis)[245] won with 32% of the 230 660 votes against the spotted lady beetle (Coleomegilla maculata lengi), the ebony jewelwing damselfly (Calopteryx maculata), a species of bumble bee (Bombus impatiens) and the six-spotted tiger beetle (Cicindela sexguttata sexguttata).
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In 1977, the Quebec Parliament declared June 24 to be Quebec's National Holiday. Historically June 24 was a holiday honouring French Canada's patron saint, St. John the Baptist, which is why it is commonly known as La Saint-Jean-Baptiste (often shortened to La St-Jean). On this day, the song "Gens du pays" by Gilles Vigneault is often heard and commonly regarded as Quebec's unofficial anthem. The festivities occur on June 23 and 24 and are celebrated all over Quebec. In cities like Québec and Montréal, great shows are organized in the main public places (such as the Abraham plains, Québec, or Maisonneuve Park, Montréal) where several of the most popular Quebec artists relay each others[clarification needed] until late at night.
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en/4843.html.txt
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Coordinates: 14h 29m 42.9487s, −62° 40′ 46.141″
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Proxima Centauri is a small, low-mass star located 4.244 light-years (1.301 pc) away from the Sun in the southern constellation of Centaurus. Its Latin name means the "nearest [star] of Centaurus". This object was discovered in 1915 by Robert Innes and is the nearest-known star to the Sun. With a quiescent apparent magnitude of 11.13, it is too faint to be seen with the naked eye. Proxima Centauri is a member of the Alpha Centauri system, being identified as component Alpha Centauri C, and is 2.18° to the southwest of the Alpha Centauri AB pair. It is currently 12,950 AU (1.94 trillion km) from AB, which it orbits with a period of about 550,000 years.
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Proxima Centauri is a red dwarf star with a mass about an eighth of the Sun's mass (M☉), and average density about 33 times that of the Sun. Because of Proxima Centauri's proximity to Earth, its angular diameter can be measured directly. Its actual diameter is about one-seventh the diameter of the Sun. Although it has a very low average luminosity, Proxima is a flare star that undergoes random dramatic increases in brightness because of magnetic activity. The star's magnetic field is created by convection throughout the stellar body, and the resulting flare activity generates a total X-ray emission similar to that produced by the Sun. The mixing of the fuel at Proxima Centauri's core through convection, and its relatively low energy-production rate, mean that it will be a main-sequence star for another four trillion years.
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In 2016, astronomers announced the discovery of Proxima Centauri b, a planet orbiting the star at a distance of roughly 0.05 AU (7.5 million km) with an orbital period of approximately 11.2 Earth days. According to updated measurements by the ESPRESSO spectrograph, its estimated mass is at least 1.17 times that of the Earth.[19] The equilibrium temperature of Proxima b is estimated to be within the range of where water could exist as liquid on its surface, thus placing it within the habitable zone of Proxima Centauri, although because Proxima Centauri is a red dwarf and a flare star, whether it could support life is disputed.
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In 1915, the Scottish astronomer Robert Innes, Director of the Union Observatory in Johannesburg, South Africa, discovered a star that had the same proper motion as Alpha Centauri.[20][21][22][23] He suggested that it be named Proxima Centauri[24] (actually Proxima Centaurus).[25] In 1917, at the Royal Observatory at the Cape of Good Hope, the Dutch astronomer Joan Voûte measured the star's trigonometric parallax at 0.755″±0.028″ and determined that Proxima Centauri was approximately the same distance from the Sun as Alpha Centauri. It was also found to be the lowest-luminosity star known at the time.[26] An equally accurate parallax determination of Proxima Centauri was made by American astronomer Harold L. Alden in 1928, who confirmed Innes's view that it is closer, with a parallax of 0.783″±0.005″.[21][24]
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In 1951, American astronomer Harlow Shapley announced that Proxima Centauri is a flare star. Examination of past photographic records showed that the star displayed a measurable increase in magnitude on about 8% of the images, making it the most active flare star then known.[27][28]
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The proximity of the star allows for detailed observation of its flare activity. In 1980, the Einstein Observatory produced a detailed X-ray energy curve of a stellar flare on Proxima Centauri. Further observations of flare activity were made with the EXOSAT and ROSAT satellites, and the X-ray emissions of smaller, solar-like flares were observed by the Japanese ASCA satellite in 1995.[29] Proxima Centauri has since been the subject of study by most X-ray observatories, including XMM-Newton and Chandra.[30]
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In 2016, the International Astronomical Union organized a Working Group on Star Names (WGSN) to catalogue and standardize proper names for stars.[31] The WGSN approved the name Proxima Centauri for this star on August 21, 2016 and it is now so included in the List of IAU approved Star Names.[32]
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Because of Proxima Centauri's southern declination, it can only be viewed south of latitude 27° N.[nb 3] Red dwarfs such as Proxima Centauri are too faint to be seen with the naked eye. Even from Alpha Centauri A or B, Proxima would only be seen as a fifth magnitude star.[33][34] It has an apparent visual magnitude of 11, so a telescope with an aperture of at least 8 cm (3.1 in) is needed to observe it, even under ideal viewing conditions—under clear, dark skies with Proxima Centauri well above the horizon.[35]
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In 2018, a superflare was observed from Proxima Centauri, the strongest flare ever seen. The optical brightness increased by a factor of 68 to approximately magnitude 6.8. It is estimated that similar flares occur around five times every year but are of such short duration, just a few minutes, that they have never been observed before.[36]
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On 2020 April 22 and 23, the New Horizons spacecraft took images of two of the nearest stars, Proxima Centauri & Wolf 359. When combined with Earth-based images, the result will be a record-setting parallax measurement.[37]
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Proxima Centauri is a red dwarf, because it belongs to the main sequence on the Hertzsprung–Russell diagram and is of spectral class M5.5. M5.5 means that it falls in the low-mass end of M-type stars.[17] Its absolute visual magnitude, or its visual magnitude as viewed from a distance of 10 parsecs (33 ly), is 15.5.[38] Its total luminosity over all wavelengths is 0.17% that of the Sun,[11] although when observed in the wavelengths of visible light the eye is most sensitive to, it is only 0.0056% as luminous as the Sun.[39] More than 85% of its radiated power is at infrared wavelengths.[40] It has a regular activity cycle of starspots.[41]
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In 2002, optical interferometry with the Very Large Telescope (VLTI) found that the angular diameter of Proxima Centauri is 1.02±0.08 mas. Because its distance is known, the actual diameter of Proxima Centauri can be calculated to be about 1/7 that of the Sun, or 1.5 times that of Jupiter. The star's mass, estimated from stellar theory, is 12.2% M☉, or 129 Jupiter masses (MJ).[42] The mass has been calculated directly, although with less precision, from observations of microlensing events to be 0.150+0.062−0.051 M☉.[43]
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Lower mass main-sequence stars have higher mean density than higher mass ones,[44] and Proxima Centauri is no exception: it has a mean density of 47.1×103 kg/m3 (47.1 g/cm3), compared with the Sun's mean density of 1.411×103 kg/m3 (1.411 g/cm3).[nb 4]
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A 1998 study of photometric variations indicates that Proxima Centauri rotates once every 83.5 days.[45] A subsequent time series analysis of chromospheric indicators in 2002 suggests a longer rotation period of 116.6±0.7 days.[46] This was subsequently ruled out in favor of a rotation period of 82.6±0.1 days.[16]
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Because of its low mass, the interior of the star is completely convective,[47] causing energy to be transferred to the exterior by the physical movement of plasma rather than through radiative processes. This convection means that the helium ash left over from the thermonuclear fusion of hydrogen does not accumulate at the core, but is instead circulated throughout the star. Unlike the Sun, which will only burn through about 10% of its total hydrogen supply before leaving the main sequence, Proxima Centauri will consume nearly all of its fuel before the fusion of hydrogen comes to an end after about 4 trillion years.[48]
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Convection is associated with the generation and persistence of a magnetic field. The magnetic energy from this field is released at the surface through stellar flares that briefly increase the overall luminosity of the star. These flares can grow as large as the star and reach temperatures measured as high as 27 million K[30]—hot enough to radiate X-rays.[49] Proxima Centauri's quiescent X-ray luminosity, approximately (4–16) × 1026 erg/s ((4–16) × 1019 W), is roughly equal to that of the much larger Sun. The peak X-ray luminosity of the largest flares can reach 1028 erg/s (1021 W).[30]
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Proxima Centauri's chromosphere is active, and its spectrum displays a strong emission line of singly ionized magnesium at a wavelength of 280 nm.[50] About 88% of the surface of Proxima Centauri may be active, a percentage that is much higher than that of the Sun even at the peak of the solar cycle. Even during quiescent periods with few or no flares, this activity increases the corona temperature of Proxima Centauri to 3.5 million K, compared to the 2 million K of the Sun's corona,[51] and its total X-ray emission is comparable to the sun's.[52] Proxima Centauri's overall activity level is considered low compared to other red dwarfs,[52] which is consistent with the star's estimated age of 4.85 × 109 years,[17] since the activity level of a red dwarf is expected to steadily wane over billions of years as its stellar rotation rate decreases.[53] The activity level also appears to vary with a period of roughly 442 days, which is shorter than the solar cycle of 11 years.[54]
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Proxima Centauri has a relatively weak stellar wind, no more than 20% of the mass loss rate of the solar wind. Because the star is much smaller than the Sun, the mass loss per unit surface area from Proxima Centauri may be eight times that from the solar surface.[55]
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A red dwarf with the mass of Proxima Centauri will remain on the main sequence for about four trillion years. As the proportion of helium increases because of hydrogen fusion, the star will become smaller and hotter, gradually transforming into a so-called "blue dwarf". Near the end of this period it will become significantly more luminous, reaching 2.5% of the Sun's luminosity (L☉) and warming up any orbiting bodies for a period of several billion years. When the hydrogen fuel is exhausted, Proxima Centauri will then evolve into a white dwarf (without passing through the red giant phase) and steadily lose any remaining heat energy.[48]
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Based on a parallax of 768.5004±0.2030 mas, published in 2018 in Gaia Data Release 2, Proxima Centauri is about 4.244 light-years (1.301 pc; 268,400 AU) from the Sun.[9] Previously published parallaxes include: 768.13±1.04 mas, in 2014 by the Research Consortium On Nearby Stars;[56] 772.33±2.42 mas, in the original Hipparcos Catalogue, in 1997;[57] 771.64±2.60 mas in the Hipparcos New Reduction, in 2007;[3] and 768.77±0.37 mas using the Hubble Space Telescope's Fine Guidance Sensors, in 1999.[10] From Earth's vantage point, Proxima is separated from Alpha Centauri by 2.18 degrees,[58] or four times the angular diameter of the full Moon.[59] Proxima also has a relatively large proper motion—moving 3.85 arcseconds per year across the sky.[60] It has a radial velocity toward the Sun of 22.2 km/s.[8]
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Among the known stars, Proxima Centauri has been the closest star to the Sun for about 32,000 years and will be so for about another 25,000 years, after which Alpha Centauri A and Alpha Centauri B will alternate approximately every 79.91 years as the closest star to the Sun. In 2001, J. García-Sánchez et al. predicted that Proxima will make its closest approach to the Sun in approximately 26,700 years, coming within 3.11 ly (0.95 pc).[61] A 2010 study by V. V. Bobylev predicted a closest approach distance of 2.90 ly (0.89 pc) in about 27,400 years,[62] followed by a 2014 study by C. A. L. Bailer-Jones predicting a perihelion approach of 3.07 ly (0.94 pc) in roughly 26,710 years.[63] Proxima Centauri is orbiting through the Milky Way at a distance from the Galactic Centre that varies from 27 to 31 kly (8.3 to 9.5 kpc), with an orbital eccentricity of 0.07.[64]
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Ever since the discovery of Proxima, it has been suspected to be a true companion of the Alpha Centauri binary star system. Data from the Hipparcos satellite, combined with ground-based observations, were consistent with the hypothesis that the three stars are a bound system. For this reason, Proxima is sometimes referred to as Alpha Centauri C. Kervella et al. (2017) used high-precision radial velocity measurements to determine with a high degree of confidence that Proxima and Alpha Centauri are gravitationally bound.[8] Proxima's orbital period around the Alpha Centauri AB barycenter is 547000+6600−4000 years with an eccentricity of 0.5±0.08; it approaches Alpha Centauri to 4300+1100−900 AU at periastron and retreats to 13000+300−100 AU at apastron.[8] At present, Proxima is 12,947 ± 260 AU (1.94 ± 0.04 trillion km) from the Alpha Centauri AB barycenter, nearly to the farthest point in its orbit.[8]
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Such a triple system can form naturally through a low-mass star being dynamically captured by a more massive binary of 1.5–2 M☉ within their embedded star cluster before the cluster disperses.[65] However, more accurate measurements of the radial velocity are needed to confirm this hypothesis.[66] If Proxima was bound to the Alpha Centauri system during its formation, the stars are likely to share the same elemental composition. The gravitational influence of Proxima might also have stirred up the Alpha Centauri protoplanetary disks. This would have increased the delivery of volatiles such as water to the dry inner regions, so possibly enriching any terrestrial planets in the system with this material.[66] Alternatively, Proxima may have been captured at a later date during an encounter, resulting in a highly eccentric orbit that was then stabilized by the galactic tide and additional stellar encounters. Such a scenario may mean that Proxima's planetary companion has had a much lower chance for orbital disruption by Alpha Centauri.[15]
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Six single stars, two binary star systems, and a triple star share a common motion through space with Proxima Centauri and the Alpha Centauri system. The space velocities of these stars are all within 10 km/s of Alpha Centauri's peculiar motion. Thus, they may form a moving group of stars, which would indicate a common point of origin,[67] such as in a star cluster.
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Proxima Centauri b, or Alpha Centauri Cb, is a planet orbiting the star at a distance of roughly 0.05 AU (7.5 million km) with an orbital period of approximately 11.2 Earth days. Its estimated mass is at least 1.3 times that of the Earth. Moreover, the equilibrium temperature of Proxima b is estimated to be within the range where water could exist as liquid on its surface; thus, placing it within the habitable zone of Proxima Centauri.[68][75][76]
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The first indications of the exoplanet Proxima Centauri b were found in 2013 by Mikko Tuomi of the University of Hertfordshire from archival observation data.[77][78] To confirm the possible discovery, a team of astronomers launched the Pale Red Dot[nb 6] project in January 2016.[79] On August 24, 2016, the team of 31 scientists from all around the world,[80] led by Guillem Anglada-Escudé of Queen Mary University of London, confirmed the existence of Proxima Centauri b[81] through a peer-reviewed article published in Nature.[68][82] The measurements were performed using two spectrographs: HARPS on the ESO 3.6 m Telescope at La Silla Observatory and UVES on the 8 m Very Large Telescope at Paranal Observatory.[68] Several attempts to detect a transit of this planet across the face of Proxima Centauri have been made. A transit-like signal appearing on September 8, 2016 was tentatively identified, using the Bright Star Survey Telescope at the Zhongshan Station in Antarctica.[83]
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Prior to this discovery, multiple measurements of the star's radial velocity constrained the maximum mass that a detectable companion to Proxima Centauri could possess.[10][84] The activity level of the star adds noise to the radial velocity measurements, complicating detection of a companion using this method.[85] In 1998, an examination of Proxima Centauri using the Faint Object Spectrograph on board the Hubble Space Telescope appeared to show evidence of a companion orbiting at a distance of about 0.5 AU.[86] A subsequent search using the Wide Field Planetary Camera 2 failed to locate any companions.[87] Astrometric measurements at the Cerro Tololo Inter-American Observatory appear to rule out a Jupiter-sized planet with an orbital period of 2−12 years.[88] A second signal in the range of 60 to 500 days was also detected, but its nature is still unclear due to stellar activity.[68]
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Proxima c is a Super-Earth about 7 times as heavy as Earth, orbiting at roughly 1.5 astronomical units (220,000,000 km) every 1,900 days (5.2 yr)[89]. Due to its large distance from Proxima Centauri, the exoplanet is unlikely to be habitable, with a low equilibrium temperature of around 39 K[90]. The planet was first reported by Italian astrophysicist Mario Damasso and his colleagues in April 2019.[90][89] Damasso's team had noticed minor movements of Proxima Centauri in the radial velocity data from the ESO's HARPS instrument, indicating a possible additional planet orbiting Proxima Centauri.[90] In 2020, the planet's existence was confirmed by precovery images from Hubble c. 1995, making it the first and closest planet ever to be directly imaged.[91] Due to the differences in how its light shines throughout its orbit, the planet is speculated to have a ring system.
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In 2017, a team of astronomers using the Atacama Large Millimeter/submillimeter Array reported detecting a belt of cold dust orbiting Proxima Centauri at a range of 1−4 AU from the star. This dust has a temperature of around 40 K and has a total estimated mass of 1% of the planet Earth. They also tentatively detected two additional features: a cold belt with a temperature of 10 K orbiting around 30 AU and a compact emission source about 1.2 arcseconds from the star. There was also a hint at an additional warm dust belt at a distance of 0.4 AU from the star.[92] However, upon further analysis, these emissions were determined to be most likely the result of a large flare emitted by the star in March 2017. The presence of dust is not needed to model the observations.[93][94]
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In 2019, a team of astronomers revisited the data from ESPRESSO about Proxima b to refine its mass. While doing so, the team found another radial velocity spike with a periodicity of 5.15 days. They estimated that if it were a planetary companion, it would be no less than 0.29 masses of the Earth.[19]
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Prior to the discovery of Proxima Centauri b, the TV documentary Alien Worlds hypothesized that a life-sustaining planet could exist in orbit around Proxima Centauri or other red dwarfs. Such a planet would lie within the habitable zone of Proxima Centauri, about 0.023–0.054 AU (3.4–8.1 million km) from the star, and would have an orbital period of 3.6–14 days.[95] A planet orbiting within this zone may experience tidal locking to the star. If the orbital eccentricity of this hypothetical planet is low, Proxima Centauri would move little in the planet's sky, and most of the surface would experience either day or night perpetually. The presence of an atmosphere could serve to redistribute the energy from the star-lit side to the far side of the planet.[96]
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Proxima Centauri's flare outbursts could erode the atmosphere of any planet in its habitable zone, but the documentary's scientists thought that this obstacle could be overcome. Gibor Basri of the University of California, Berkeley, mentioned that "no one [has] found any showstoppers to habitability". For example, one concern was that the torrents of charged particles from the star's flares could strip the atmosphere off any nearby planet. If the planet had a strong magnetic field, the field would deflect the particles from the atmosphere; even the slow rotation of a tidally locked planet that spins once for every time it orbits its star would be enough to generate a magnetic field, as long as part of the planet's interior remained molten.[97]
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Other scientists, especially proponents of the rare-Earth hypothesis,[98] disagree that red dwarfs can sustain life. Any exoplanet in this star's habitable zone would likely be tidally locked, resulting in a relatively weak planetary magnetic moment, leading to strong atmospheric erosion by coronal mass ejections from Proxima Centauri.[99]
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Because of the star's proximity to Earth, Proxima Centauri has been proposed as a flyby destination for interstellar travel.[100]
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Proxima currently moves toward Earth at a rate of 22.2 km/s.[8] After 26,700 years, when it will come within 3.11 light-years, it will begin to move farther away.[61]
|
70 |
+
|
71 |
+
If non-nuclear, conventional propulsion technologies are used, the flight of a spacecraft to a planet orbiting Proxima Centauri would probably require thousands of years.[101] For example, Voyager 1, which is now travelling 17 km/s (38,000 mph)[102] relative to the Sun, would reach Proxima in 73,775 years, were the spacecraft travelling in the direction of that star. A slow-moving probe would have only several tens of thousands of years to catch Proxima Centauri near its closest approach, and could end up watching it recede into the distance.[103]
|
72 |
+
|
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+
Nuclear pulse propulsion might enable such interstellar travel with a trip timescale of a century, inspiring several studies such as Project Orion, Project Daedalus, and Project Longshot.[103]
|
74 |
+
|
75 |
+
Project Breakthrough Starshot aims to reach the Alpha Centauri system within the first half of the 21st century, with microprobes travelling at 20% of the speed of light propelled by around 100 gigawatts of Earth-based lasers.[104] The probes would perform a fly-by of Proxima Centauri to take photos and collect data of its planets' atmospheric compositions. It would take 4.22 years for the information collected to be sent back to Earth.[105]
|
76 |
+
|
77 |
+
From Proxima Centauri, the Sun would appear as a bright 0.4-magnitude star in the constellation Cassiopeia, similar to that of Achernar from Earth.[nb 7]
|
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+
|
79 |
+
where
|
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+
|
88 |
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|
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+
ρ
|
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+
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+
⊙
|
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+
|
93 |
+
|
94 |
+
|
95 |
+
|
96 |
+
|
97 |
+
|
98 |
+
|
99 |
+
|
100 |
+
|
101 |
+
{\displaystyle {\begin{smallmatrix}\rho _{\odot }\end{smallmatrix}}}
|
102 |
+
|
103 |
+
is the average solar density.
|
104 |
+
See:
|
en/4844.html.txt
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1 |
+
Coordinates: 14h 29m 42.9487s, −62° 40′ 46.141″
|
2 |
+
|
3 |
+
Proxima Centauri is a small, low-mass star located 4.244 light-years (1.301 pc) away from the Sun in the southern constellation of Centaurus. Its Latin name means the "nearest [star] of Centaurus". This object was discovered in 1915 by Robert Innes and is the nearest-known star to the Sun. With a quiescent apparent magnitude of 11.13, it is too faint to be seen with the naked eye. Proxima Centauri is a member of the Alpha Centauri system, being identified as component Alpha Centauri C, and is 2.18° to the southwest of the Alpha Centauri AB pair. It is currently 12,950 AU (1.94 trillion km) from AB, which it orbits with a period of about 550,000 years.
|
4 |
+
|
5 |
+
Proxima Centauri is a red dwarf star with a mass about an eighth of the Sun's mass (M☉), and average density about 33 times that of the Sun. Because of Proxima Centauri's proximity to Earth, its angular diameter can be measured directly. Its actual diameter is about one-seventh the diameter of the Sun. Although it has a very low average luminosity, Proxima is a flare star that undergoes random dramatic increases in brightness because of magnetic activity. The star's magnetic field is created by convection throughout the stellar body, and the resulting flare activity generates a total X-ray emission similar to that produced by the Sun. The mixing of the fuel at Proxima Centauri's core through convection, and its relatively low energy-production rate, mean that it will be a main-sequence star for another four trillion years.
|
6 |
+
|
7 |
+
In 2016, astronomers announced the discovery of Proxima Centauri b, a planet orbiting the star at a distance of roughly 0.05 AU (7.5 million km) with an orbital period of approximately 11.2 Earth days. According to updated measurements by the ESPRESSO spectrograph, its estimated mass is at least 1.17 times that of the Earth.[19] The equilibrium temperature of Proxima b is estimated to be within the range of where water could exist as liquid on its surface, thus placing it within the habitable zone of Proxima Centauri, although because Proxima Centauri is a red dwarf and a flare star, whether it could support life is disputed.
|
8 |
+
|
9 |
+
In 1915, the Scottish astronomer Robert Innes, Director of the Union Observatory in Johannesburg, South Africa, discovered a star that had the same proper motion as Alpha Centauri.[20][21][22][23] He suggested that it be named Proxima Centauri[24] (actually Proxima Centaurus).[25] In 1917, at the Royal Observatory at the Cape of Good Hope, the Dutch astronomer Joan Voûte measured the star's trigonometric parallax at 0.755″±0.028″ and determined that Proxima Centauri was approximately the same distance from the Sun as Alpha Centauri. It was also found to be the lowest-luminosity star known at the time.[26] An equally accurate parallax determination of Proxima Centauri was made by American astronomer Harold L. Alden in 1928, who confirmed Innes's view that it is closer, with a parallax of 0.783″±0.005″.[21][24]
|
10 |
+
|
11 |
+
In 1951, American astronomer Harlow Shapley announced that Proxima Centauri is a flare star. Examination of past photographic records showed that the star displayed a measurable increase in magnitude on about 8% of the images, making it the most active flare star then known.[27][28]
|
12 |
+
The proximity of the star allows for detailed observation of its flare activity. In 1980, the Einstein Observatory produced a detailed X-ray energy curve of a stellar flare on Proxima Centauri. Further observations of flare activity were made with the EXOSAT and ROSAT satellites, and the X-ray emissions of smaller, solar-like flares were observed by the Japanese ASCA satellite in 1995.[29] Proxima Centauri has since been the subject of study by most X-ray observatories, including XMM-Newton and Chandra.[30]
|
13 |
+
|
14 |
+
In 2016, the International Astronomical Union organized a Working Group on Star Names (WGSN) to catalogue and standardize proper names for stars.[31] The WGSN approved the name Proxima Centauri for this star on August 21, 2016 and it is now so included in the List of IAU approved Star Names.[32]
|
15 |
+
|
16 |
+
Because of Proxima Centauri's southern declination, it can only be viewed south of latitude 27° N.[nb 3] Red dwarfs such as Proxima Centauri are too faint to be seen with the naked eye. Even from Alpha Centauri A or B, Proxima would only be seen as a fifth magnitude star.[33][34] It has an apparent visual magnitude of 11, so a telescope with an aperture of at least 8 cm (3.1 in) is needed to observe it, even under ideal viewing conditions—under clear, dark skies with Proxima Centauri well above the horizon.[35]
|
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+
|
18 |
+
In 2018, a superflare was observed from Proxima Centauri, the strongest flare ever seen. The optical brightness increased by a factor of 68 to approximately magnitude 6.8. It is estimated that similar flares occur around five times every year but are of such short duration, just a few minutes, that they have never been observed before.[36]
|
19 |
+
|
20 |
+
On 2020 April 22 and 23, the New Horizons spacecraft took images of two of the nearest stars, Proxima Centauri & Wolf 359. When combined with Earth-based images, the result will be a record-setting parallax measurement.[37]
|
21 |
+
|
22 |
+
Proxima Centauri is a red dwarf, because it belongs to the main sequence on the Hertzsprung–Russell diagram and is of spectral class M5.5. M5.5 means that it falls in the low-mass end of M-type stars.[17] Its absolute visual magnitude, or its visual magnitude as viewed from a distance of 10 parsecs (33 ly), is 15.5.[38] Its total luminosity over all wavelengths is 0.17% that of the Sun,[11] although when observed in the wavelengths of visible light the eye is most sensitive to, it is only 0.0056% as luminous as the Sun.[39] More than 85% of its radiated power is at infrared wavelengths.[40] It has a regular activity cycle of starspots.[41]
|
23 |
+
|
24 |
+
In 2002, optical interferometry with the Very Large Telescope (VLTI) found that the angular diameter of Proxima Centauri is 1.02±0.08 mas. Because its distance is known, the actual diameter of Proxima Centauri can be calculated to be about 1/7 that of the Sun, or 1.5 times that of Jupiter. The star's mass, estimated from stellar theory, is 12.2% M☉, or 129 Jupiter masses (MJ).[42] The mass has been calculated directly, although with less precision, from observations of microlensing events to be 0.150+0.062−0.051 M☉.[43]
|
25 |
+
|
26 |
+
Lower mass main-sequence stars have higher mean density than higher mass ones,[44] and Proxima Centauri is no exception: it has a mean density of 47.1×103 kg/m3 (47.1 g/cm3), compared with the Sun's mean density of 1.411×103 kg/m3 (1.411 g/cm3).[nb 4]
|
27 |
+
|
28 |
+
A 1998 study of photometric variations indicates that Proxima Centauri rotates once every 83.5 days.[45] A subsequent time series analysis of chromospheric indicators in 2002 suggests a longer rotation period of 116.6±0.7 days.[46] This was subsequently ruled out in favor of a rotation period of 82.6±0.1 days.[16]
|
29 |
+
|
30 |
+
Because of its low mass, the interior of the star is completely convective,[47] causing energy to be transferred to the exterior by the physical movement of plasma rather than through radiative processes. This convection means that the helium ash left over from the thermonuclear fusion of hydrogen does not accumulate at the core, but is instead circulated throughout the star. Unlike the Sun, which will only burn through about 10% of its total hydrogen supply before leaving the main sequence, Proxima Centauri will consume nearly all of its fuel before the fusion of hydrogen comes to an end after about 4 trillion years.[48]
|
31 |
+
|
32 |
+
Convection is associated with the generation and persistence of a magnetic field. The magnetic energy from this field is released at the surface through stellar flares that briefly increase the overall luminosity of the star. These flares can grow as large as the star and reach temperatures measured as high as 27 million K[30]—hot enough to radiate X-rays.[49] Proxima Centauri's quiescent X-ray luminosity, approximately (4–16) × 1026 erg/s ((4–16) × 1019 W), is roughly equal to that of the much larger Sun. The peak X-ray luminosity of the largest flares can reach 1028 erg/s (1021 W).[30]
|
33 |
+
|
34 |
+
Proxima Centauri's chromosphere is active, and its spectrum displays a strong emission line of singly ionized magnesium at a wavelength of 280 nm.[50] About 88% of the surface of Proxima Centauri may be active, a percentage that is much higher than that of the Sun even at the peak of the solar cycle. Even during quiescent periods with few or no flares, this activity increases the corona temperature of Proxima Centauri to 3.5 million K, compared to the 2 million K of the Sun's corona,[51] and its total X-ray emission is comparable to the sun's.[52] Proxima Centauri's overall activity level is considered low compared to other red dwarfs,[52] which is consistent with the star's estimated age of 4.85 × 109 years,[17] since the activity level of a red dwarf is expected to steadily wane over billions of years as its stellar rotation rate decreases.[53] The activity level also appears to vary with a period of roughly 442 days, which is shorter than the solar cycle of 11 years.[54]
|
35 |
+
|
36 |
+
Proxima Centauri has a relatively weak stellar wind, no more than 20% of the mass loss rate of the solar wind. Because the star is much smaller than the Sun, the mass loss per unit surface area from Proxima Centauri may be eight times that from the solar surface.[55]
|
37 |
+
|
38 |
+
A red dwarf with the mass of Proxima Centauri will remain on the main sequence for about four trillion years. As the proportion of helium increases because of hydrogen fusion, the star will become smaller and hotter, gradually transforming into a so-called "blue dwarf". Near the end of this period it will become significantly more luminous, reaching 2.5% of the Sun's luminosity (L☉) and warming up any orbiting bodies for a period of several billion years. When the hydrogen fuel is exhausted, Proxima Centauri will then evolve into a white dwarf (without passing through the red giant phase) and steadily lose any remaining heat energy.[48]
|
39 |
+
|
40 |
+
Based on a parallax of 768.5004±0.2030 mas, published in 2018 in Gaia Data Release 2, Proxima Centauri is about 4.244 light-years (1.301 pc; 268,400 AU) from the Sun.[9] Previously published parallaxes include: 768.13±1.04 mas, in 2014 by the Research Consortium On Nearby Stars;[56] 772.33±2.42 mas, in the original Hipparcos Catalogue, in 1997;[57] 771.64±2.60 mas in the Hipparcos New Reduction, in 2007;[3] and 768.77±0.37 mas using the Hubble Space Telescope's Fine Guidance Sensors, in 1999.[10] From Earth's vantage point, Proxima is separated from Alpha Centauri by 2.18 degrees,[58] or four times the angular diameter of the full Moon.[59] Proxima also has a relatively large proper motion—moving 3.85 arcseconds per year across the sky.[60] It has a radial velocity toward the Sun of 22.2 km/s.[8]
|
41 |
+
|
42 |
+
Among the known stars, Proxima Centauri has been the closest star to the Sun for about 32,000 years and will be so for about another 25,000 years, after which Alpha Centauri A and Alpha Centauri B will alternate approximately every 79.91 years as the closest star to the Sun. In 2001, J. García-Sánchez et al. predicted that Proxima will make its closest approach to the Sun in approximately 26,700 years, coming within 3.11 ly (0.95 pc).[61] A 2010 study by V. V. Bobylev predicted a closest approach distance of 2.90 ly (0.89 pc) in about 27,400 years,[62] followed by a 2014 study by C. A. L. Bailer-Jones predicting a perihelion approach of 3.07 ly (0.94 pc) in roughly 26,710 years.[63] Proxima Centauri is orbiting through the Milky Way at a distance from the Galactic Centre that varies from 27 to 31 kly (8.3 to 9.5 kpc), with an orbital eccentricity of 0.07.[64]
|
43 |
+
|
44 |
+
Ever since the discovery of Proxima, it has been suspected to be a true companion of the Alpha Centauri binary star system. Data from the Hipparcos satellite, combined with ground-based observations, were consistent with the hypothesis that the three stars are a bound system. For this reason, Proxima is sometimes referred to as Alpha Centauri C. Kervella et al. (2017) used high-precision radial velocity measurements to determine with a high degree of confidence that Proxima and Alpha Centauri are gravitationally bound.[8] Proxima's orbital period around the Alpha Centauri AB barycenter is 547000+6600−4000 years with an eccentricity of 0.5±0.08; it approaches Alpha Centauri to 4300+1100−900 AU at periastron and retreats to 13000+300−100 AU at apastron.[8] At present, Proxima is 12,947 ± 260 AU (1.94 ± 0.04 trillion km) from the Alpha Centauri AB barycenter, nearly to the farthest point in its orbit.[8]
|
45 |
+
|
46 |
+
Such a triple system can form naturally through a low-mass star being dynamically captured by a more massive binary of 1.5–2 M☉ within their embedded star cluster before the cluster disperses.[65] However, more accurate measurements of the radial velocity are needed to confirm this hypothesis.[66] If Proxima was bound to the Alpha Centauri system during its formation, the stars are likely to share the same elemental composition. The gravitational influence of Proxima might also have stirred up the Alpha Centauri protoplanetary disks. This would have increased the delivery of volatiles such as water to the dry inner regions, so possibly enriching any terrestrial planets in the system with this material.[66] Alternatively, Proxima may have been captured at a later date during an encounter, resulting in a highly eccentric orbit that was then stabilized by the galactic tide and additional stellar encounters. Such a scenario may mean that Proxima's planetary companion has had a much lower chance for orbital disruption by Alpha Centauri.[15]
|
47 |
+
|
48 |
+
Six single stars, two binary star systems, and a triple star share a common motion through space with Proxima Centauri and the Alpha Centauri system. The space velocities of these stars are all within 10 km/s of Alpha Centauri's peculiar motion. Thus, they may form a moving group of stars, which would indicate a common point of origin,[67] such as in a star cluster.
|
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+
|
50 |
+
Proxima Centauri b, or Alpha Centauri Cb, is a planet orbiting the star at a distance of roughly 0.05 AU (7.5 million km) with an orbital period of approximately 11.2 Earth days. Its estimated mass is at least 1.3 times that of the Earth. Moreover, the equilibrium temperature of Proxima b is estimated to be within the range where water could exist as liquid on its surface; thus, placing it within the habitable zone of Proxima Centauri.[68][75][76]
|
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+
|
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+
The first indications of the exoplanet Proxima Centauri b were found in 2013 by Mikko Tuomi of the University of Hertfordshire from archival observation data.[77][78] To confirm the possible discovery, a team of astronomers launched the Pale Red Dot[nb 6] project in January 2016.[79] On August 24, 2016, the team of 31 scientists from all around the world,[80] led by Guillem Anglada-Escudé of Queen Mary University of London, confirmed the existence of Proxima Centauri b[81] through a peer-reviewed article published in Nature.[68][82] The measurements were performed using two spectrographs: HARPS on the ESO 3.6 m Telescope at La Silla Observatory and UVES on the 8 m Very Large Telescope at Paranal Observatory.[68] Several attempts to detect a transit of this planet across the face of Proxima Centauri have been made. A transit-like signal appearing on September 8, 2016 was tentatively identified, using the Bright Star Survey Telescope at the Zhongshan Station in Antarctica.[83]
|
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+
|
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+
Prior to this discovery, multiple measurements of the star's radial velocity constrained the maximum mass that a detectable companion to Proxima Centauri could possess.[10][84] The activity level of the star adds noise to the radial velocity measurements, complicating detection of a companion using this method.[85] In 1998, an examination of Proxima Centauri using the Faint Object Spectrograph on board the Hubble Space Telescope appeared to show evidence of a companion orbiting at a distance of about 0.5 AU.[86] A subsequent search using the Wide Field Planetary Camera 2 failed to locate any companions.[87] Astrometric measurements at the Cerro Tololo Inter-American Observatory appear to rule out a Jupiter-sized planet with an orbital period of 2−12 years.[88] A second signal in the range of 60 to 500 days was also detected, but its nature is still unclear due to stellar activity.[68]
|
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Proxima c is a Super-Earth about 7 times as heavy as Earth, orbiting at roughly 1.5 astronomical units (220,000,000 km) every 1,900 days (5.2 yr)[89]. Due to its large distance from Proxima Centauri, the exoplanet is unlikely to be habitable, with a low equilibrium temperature of around 39 K[90]. The planet was first reported by Italian astrophysicist Mario Damasso and his colleagues in April 2019.[90][89] Damasso's team had noticed minor movements of Proxima Centauri in the radial velocity data from the ESO's HARPS instrument, indicating a possible additional planet orbiting Proxima Centauri.[90] In 2020, the planet's existence was confirmed by precovery images from Hubble c. 1995, making it the first and closest planet ever to be directly imaged.[91] Due to the differences in how its light shines throughout its orbit, the planet is speculated to have a ring system.
|
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+
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+
In 2017, a team of astronomers using the Atacama Large Millimeter/submillimeter Array reported detecting a belt of cold dust orbiting Proxima Centauri at a range of 1−4 AU from the star. This dust has a temperature of around 40 K and has a total estimated mass of 1% of the planet Earth. They also tentatively detected two additional features: a cold belt with a temperature of 10 K orbiting around 30 AU and a compact emission source about 1.2 arcseconds from the star. There was also a hint at an additional warm dust belt at a distance of 0.4 AU from the star.[92] However, upon further analysis, these emissions were determined to be most likely the result of a large flare emitted by the star in March 2017. The presence of dust is not needed to model the observations.[93][94]
|
59 |
+
|
60 |
+
In 2019, a team of astronomers revisited the data from ESPRESSO about Proxima b to refine its mass. While doing so, the team found another radial velocity spike with a periodicity of 5.15 days. They estimated that if it were a planetary companion, it would be no less than 0.29 masses of the Earth.[19]
|
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+
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+
Prior to the discovery of Proxima Centauri b, the TV documentary Alien Worlds hypothesized that a life-sustaining planet could exist in orbit around Proxima Centauri or other red dwarfs. Such a planet would lie within the habitable zone of Proxima Centauri, about 0.023–0.054 AU (3.4–8.1 million km) from the star, and would have an orbital period of 3.6–14 days.[95] A planet orbiting within this zone may experience tidal locking to the star. If the orbital eccentricity of this hypothetical planet is low, Proxima Centauri would move little in the planet's sky, and most of the surface would experience either day or night perpetually. The presence of an atmosphere could serve to redistribute the energy from the star-lit side to the far side of the planet.[96]
|
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Proxima Centauri's flare outbursts could erode the atmosphere of any planet in its habitable zone, but the documentary's scientists thought that this obstacle could be overcome. Gibor Basri of the University of California, Berkeley, mentioned that "no one [has] found any showstoppers to habitability". For example, one concern was that the torrents of charged particles from the star's flares could strip the atmosphere off any nearby planet. If the planet had a strong magnetic field, the field would deflect the particles from the atmosphere; even the slow rotation of a tidally locked planet that spins once for every time it orbits its star would be enough to generate a magnetic field, as long as part of the planet's interior remained molten.[97]
|
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+
|
66 |
+
Other scientists, especially proponents of the rare-Earth hypothesis,[98] disagree that red dwarfs can sustain life. Any exoplanet in this star's habitable zone would likely be tidally locked, resulting in a relatively weak planetary magnetic moment, leading to strong atmospheric erosion by coronal mass ejections from Proxima Centauri.[99]
|
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+
|
68 |
+
Because of the star's proximity to Earth, Proxima Centauri has been proposed as a flyby destination for interstellar travel.[100]
|
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+
Proxima currently moves toward Earth at a rate of 22.2 km/s.[8] After 26,700 years, when it will come within 3.11 light-years, it will begin to move farther away.[61]
|
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If non-nuclear, conventional propulsion technologies are used, the flight of a spacecraft to a planet orbiting Proxima Centauri would probably require thousands of years.[101] For example, Voyager 1, which is now travelling 17 km/s (38,000 mph)[102] relative to the Sun, would reach Proxima in 73,775 years, were the spacecraft travelling in the direction of that star. A slow-moving probe would have only several tens of thousands of years to catch Proxima Centauri near its closest approach, and could end up watching it recede into the distance.[103]
|
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+
|
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+
Nuclear pulse propulsion might enable such interstellar travel with a trip timescale of a century, inspiring several studies such as Project Orion, Project Daedalus, and Project Longshot.[103]
|
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+
|
75 |
+
Project Breakthrough Starshot aims to reach the Alpha Centauri system within the first half of the 21st century, with microprobes travelling at 20% of the speed of light propelled by around 100 gigawatts of Earth-based lasers.[104] The probes would perform a fly-by of Proxima Centauri to take photos and collect data of its planets' atmospheric compositions. It would take 4.22 years for the information collected to be sent back to Earth.[105]
|
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+
|
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+
From Proxima Centauri, the Sun would appear as a bright 0.4-magnitude star in the constellation Cassiopeia, similar to that of Achernar from Earth.[nb 7]
|
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|
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where
|
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ρ
|
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⊙
|
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|
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|
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|
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+
{\displaystyle {\begin{smallmatrix}\rho _{\odot }\end{smallmatrix}}}
|
102 |
+
|
103 |
+
is the average solar density.
|
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See:
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See text
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A plum is a fruit of the subgenus Prunus of the genus Prunus. The subgenus is distinguished from other subgenera (peaches, cherries, bird cherries, etc.) in the shoots having terminal bud and solitary side buds (not clustered), the flowers in groups of one to five together on short stems, and the fruit having a groove running down one side and a smooth stone (or pit).
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Mature plum fruit may have a dusty-white waxy coating that gives them a glaucous appearance. This is an epicuticular wax coating and is known as "wax bloom". Dried plum fruits are called "dried plums" or prunes, although, in many countries, prunes are a distinct type of dried plum having a wrinkled appearance (Li hing mui for instance).[1]
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Plums may have been one of the first fruits domesticated by humans.[2] Three of the most abundant cultivars are not found in the wild, only around human settlements: Prunus domestica has been traced to East European and Caucasian mountains, while Prunus salicina and Prunus simonii originated in Asia. Plum remains have been found in Neolithic age archaeological sites along with olives, grapes and figs.[3][4] According to Ken Albala, plums originated in Iran.[5]
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The name plum derived from Old English plume or "plum, plum tree," which extended from Germanic language or Middle Dutch, and Latin prūnum, from Ancient Greek προῦμνον (proumnon), believed to be a loanword from Asia Minor.[1][6] In the late 18th century, the word, plum, was used to indicate "something desirable", probably in reference to tasty fruit pieces in desserts.[6]
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Plums are a diverse group of species. The commercially important plum trees are medium-sized, usually pruned to 5–6 metres height. The tree is of medium hardiness.[7] Without pruning, the trees can reach 12 metres in height and spread across 10 metres. They blossom in different months in different parts of the world; for example, in about January in Taiwan and early April in the United Kingdom.[8]
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Fruits are usually of medium size, between 2 and 7 centimetres in diameter, globose to oval. The flesh is firm and juicy. The fruit's peel is smooth, with a natural waxy surface that adheres to the flesh. The plum is a drupe, meaning its fleshy fruit surrounds a single hard seed.
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Plum cultivars include:
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Damsons
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Greengages
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Mirabelles
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Victoria plums
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When it flowers in the early spring, a plum tree will be covered in blossoms, and in a good year approximately 50% of the flowers will be pollinated and become plums. Flowering starts after 80 growing degree days.
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If the weather is too dry, the plums will not develop past a certain stage, but will fall from the tree while still tiny, green buds, and if it is unseasonably wet or if the plums are not harvested as soon as they are ripe, the fruit may develop a fungal condition called brown rot. Brown rot is not toxic, and some affected areas can be cut out of the fruit, but unless the rot is caught immediately, the fruit will no longer be edible. Plum is used as a food plant by the larvae of some Lepidoptera, including November moth, willow beauty and short-cloaked moth.
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The taste of the plum fruit ranges from sweet to tart; the skin itself may be particularly tart. It is juicy and can be eaten fresh or used in jam-making or other recipes. Plum juice can be fermented into plum wine. In central England, a cider-like alcoholic beverage known as plum jerkum is made from plums. Dried, salted plums are used as a snack, sometimes known as saladito or salao. Various flavors of dried plum are available at Chinese grocers and specialty stores worldwide. They tend to be much drier than the standard prune. Cream, ginseng, spicy, and salty are among the common varieties. Licorice is generally used to intensify the flavor of these plums and is used to make salty plum drinks and toppings for shaved ice or baobing. Pickled plums are another type of preserve available in Asia and international specialty stores. The Japanese variety, called umeboshi, is often used for rice balls, called onigiri or omusubi. The ume, from which umeboshi are made, is more closely related, however, to the apricot than to the plum. In the Balkans, plum is converted into an alcoholic drink named slivovitz (plum brandy) (Serbian: šljivovica).[9][10] A large number of plums, of the Damson variety, are also grown in Hungary, where they are called szilva and are used to make lekvar (a plum paste jam), palinka (traditional fruit brandy), plum dumplings, and other foods. In Romania, 80% of the plum production is used to create a similar brandy, called țuică.[11]
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As with many other members of the rose family, plum kernels contain cyanogenic glycosides, including amygdalin.[12] Prune kernel oil is made from the fleshy inner part of the pit of the plum. Though not available commercially, the wood of plum trees is used by hobbyists and other private woodworkers for musical instruments, knife handles, inlays, and similar small projects.[13]
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Plums have many species, and taxonomists differ on the count. Depending on the taxonomist, between 19 and 40 species of plum exist. From this diversity only two species, the hexaploid European plum (Prunus domestica) and the diploid Japanese plum (Prunus salicina and hybrids), are of worldwide commercial significance. The origin of these commercially important species is uncertain but may have involved P. cerasifera and possibly P. spinosa as ancestors. Other species of plum variously originated in Europe, Asia and America.[15]
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The subgenus Prunus is divided into three sections:
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Sect. Prunus (Old World plums) – leaves in bud rolled inwards; flowers 1-3 together; fruit smooth, often wax-bloomed
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Sect. Prunocerasus (New World plums) – leaves in bud folded inwards; flowers 3–5 together; fruit smooth, often wax-bloomed
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Sect. Armeniaca (apricots) – leaves in bud rolled inwards; flowers very short-stalked; fruit velvety; treated as a distinct subgenus by some authors
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In certain parts of the world, some fruits are called plums and are quite different from fruits known as plums in Europe or the Americas. For example, marian plums are popular in Thailand, Malaysia and Indonesia, otherwise also known as gandaria, plum mango, ma-praang, ma-yong, ramania, kundang, rembunia or setar.[17] Another example is the loquat, also known as Japanese plum and Japanese medlar, as well as nispero, bibassier and wollmispel elsewhere.[18][19] In South Asia and Southeast Asia, Jambul, a fruit from tropical tree in family Myrtaceae, is similarly sometimes referred to 'damson plums', and it is different from damson plums found in Europe and Americas.[20] Jambul is also called as Java plum, Malabar plum, Jaman, Jamun, Jamblang, Jiwat, Salam, Duhat, Koeli, Jambuláo or Koriang.
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In 2016, global production of plums and sloes was 12.1 million tonnes, led by China with 55% of the world total (table).[14] Other major producers were Romania, Serbia, and the United States (table).
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Raw plums are 87% water, 11% carbohydrates, 1% protein, and less than 1% fat (table). In a 100-gram (3 1⁄2-ounce) reference serving, raw plums supply 192 kilojoules (46 kilocalories) of food energy and are a moderate source only of vitamin C (12% Daily Value), with no other nutrients in significant content (table).
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The PlayStation 2 (officially branded as PS2) is a home video game console developed and marketed by Sony Computer Entertainment. It was first released in Japan on March 4, 2000, in North America on October 26, 2000, and in Europe and Australia on November 24, 2000, and is the successor to the original PlayStation, as well as the second installment in the PlayStation console line-up. A sixth-generation console, it competed with Sega's Dreamcast, Nintendo's GameCube, and Microsoft's original Xbox.
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Announced in 1999, the PS2 offered backward-compatibility for its predecessor's DualShock controller, as well as its games. The PS2 is the best-selling video game console of all time, having sold over 155 million units worldwide, as confirmed by Sony.[13] Over 3,800 game titles have been released for the PS2, with over 1.5 billion copies sold.[14] Sony later manufactured several smaller, lighter revisions of the console known as Slimline models in 2004.
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Even with the release of its successor, the PlayStation 3, the PS2 remained popular well into the seventh generation, and continued to be produced until 2013, when Sony finally announced it had been discontinued after over twelve years of production – one of the longest lifespans of a video game console. Despite the announcement, new games for the console continued to be produced until the end of 2013, including Final Fantasy XI: Seekers of Adoulin for Japan, FIFA 13 for North America, and Pro Evolution Soccer 2014 for Europe. Repair services for the system in Japan ended on September 7, 2018.
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Though Sony has kept details of the PlayStation 2's development secret, work on the console began around the time that the original PlayStation was released (in late 1994).[15] Insiders stated that it was developed in the U.S. West Coast by former members of Argonaut Software.[16] By 1997 word had leaked to the press that the console would have backward-compatibility with the original PlayStation, a built-in DVD player, and Internet connectivity.[16][17] Sony announced the PlayStation 2 (PS2) on March 1, 1999. The video game console was positioned as a competitor to Sega's Dreamcast, the first sixth-generation console to be released, although ultimately the main rivals of the PS2 were Nintendo's GameCube and Microsoft's Xbox.[18][19] The Dreamcast itself launched very successfully in North America later that year, selling over 500,000 units within two weeks.[20]
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Soon after the Dreamcast's North American launch, Sony unveiled the PlayStation 2 at the Tokyo Game Show on September 20, 1999.[21] Sony showed fully playable demos of upcoming PlayStation 2 games including Gran Turismo 2000 (later released as Gran Turismo 3: A-Spec) and Tekken Tag Tournament – which showed the console's graphic abilities and power.[22]
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The PS2 was launched in March 2000 in Japan, October in North America, and November in Europe. Sales of the console, games and accessories pulled in $250 million on the first day, beating the $97 million made on the first day of the Dreamcast.[23] Directly after its release, it was difficult to find PS2 units on retailer shelves[24] due to manufacturing delays.[25] Another option was purchasing the console online through auction websites such as eBay, where people paid over a thousand dollars for the console.[26] The PS2 initially sold well partly on the basis of the strength of the PlayStation brand and the console's backward-compatibility, selling over 980,000 units in Japan by March 5, 2000, one day after launch.[27] This allowed the PS2 to tap the large install base established by the PlayStation – another major selling point over the competition. Later, Sony added new development kits for game developers and more PS2 units for consumers. The PS2's built-in functionality also expanded its audience beyond the gamer,[5] as its debut pricing was the same or less than a standalone DVD player. This made the console a low cost entry into the home theater market.[28]
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The success of the PS2 at the end of 2000 caused Sega problems both financially and competitively, and Sega announced the discontinuation of the Dreamcast in March 2001, just 18 months after its successful Western launch. Despite the Dreamcast still receiving support through 2001, the PS2 remained the only sixth generation console for over 6 months before it faced competition from new rivals: Nintendo's GameCube and Microsoft's Xbox. Many analysts predicted a close three-way matchup among the three consoles. The Xbox had the most powerful hardware, while the GameCube was the least expensive console, and Nintendo changed its policy to encourage third-party developers. While the PlayStation 2 theoretically had the weakest specification of the three, it had a head start due to its installed base plus strong developer commitment, as well as a built-in DVD player (the Xbox required an adapter, while the GameCube lacked support entirely).[29] While the PlayStation 2's initial games lineup was considered mediocre, this changed during the 2001 holiday season with the release of several blockbuster games that maintained the PS2's sales momentum and held off its newer rivals. Sony also countered the Xbox by temporarily securing PlayStation 2 exclusives for highly anticipated games such as the Grand Theft Auto series and Metal Gear Solid 2: Sons of Liberty.[30]
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Sony cut the price of the console in May 2002 from US$299 to $199 in North America,[31] making it the same price as the GameCube and $100 less than the Xbox. It also planned to cut the price in Japan around that time.[32] It cut the price twice in Japan in 2003.[33] In 2006, Sony cut the cost of the console in anticipation of the release of the PlayStation 3.[33]
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Sony, unlike Sega with its Dreamcast, originally placed little emphasis on online gaming during its first few years, although that changed upon the launch of the online-capable Xbox. Coinciding with the release of Xbox Live, Sony released the PlayStation Network Adapter in late 2002, with several online first–party titles released alongside it, such as SOCOM: U.S. Navy SEALs to demonstrate its active support for Internet play.[34] Sony also advertised heavily, and its online model had the support of Electronic Arts (EA); EA did not offer online Xbox titles until 2004. Although Sony and Nintendo both started out late, and although both followed a decentralized model of online gaming where the responsibility is up to the developer to provide the servers, Sony's moves made online gaming a major selling point of the PS2.
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In September 2004, in time for the launch of Grand Theft Auto: San Andreas, Sony revealed a newer, slimmer PS2. In preparation for the launch of the new models (SCPH-700xx-9000x), Sony stopped making the older models (SCPH-3000x-500xx) to let the distribution channel empty its stock of the units.[citation needed] After an apparent manufacturing issue – Sony reportedly underestimated demand – caused some initial slowdown in producing the new unit caused in part by shortages between the time the old units were cleared out and the new units were ready. The issue was compounded in Britain when a Russian oil tanker became stuck in the Suez Canal, blocking a ship from China carrying PS2s bound for the UK. During one week in November, British sales totalled 6,000 units – compared to 70,000 units a few weeks prior.[35] There were shortages in more than 1,700 stores in North America on the day before Christmas.[36]
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Software for the PlayStation 2 was distributed primarily on DVD-ROM, with some titles being published on CD-ROM. In addition, the console can play audio CDs and DVD movies and is backward-compatible with almost all original PlayStation games. The PlayStation 2 also supports PlayStation memory cards and controllers, although original PlayStation memory cards will only work with original PlayStation games[37] and the controllers may not support all functions (such as analog buttons) for PlayStation 2 games.
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The standard PlayStation 2 memory card has an 8 MB capacity.[38] There are a variety of non-Sony manufactured memory cards available for the PlayStation 2, allowing for a memory capacity larger than the standard 8 MB.
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The console also features 2 USB ports, and 1 IEEE 1394 (Firewire) port (SCPH-10000 to 3900x only). A hard disk drive can be installed in an expansion bay on the back of the console, and is required to play certain games, notably the popular Final Fantasy XI.[39] This is only available on certain models.
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The console uses the Emotion Engine CPU, custom-designed by Sony and Toshiba and based on the MIPS architecture with a floating point performance of 6.2 GFLOPS.[40] The GPU is likewise custom-designed for the console and called the Graphics Synthesizer, with a fillrate of 2.4 gigapixels/second, capable of rendering up to 75 million polygons per second.[41] When accounting for features such as lighting, texture mapping, artificial intelligence, and game physics, it has a real-world performance of 3 million to 16 million polygons per second.[41][42]
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The PlayStation 2 may natively output video resolutions on SDTV and HDTV from 480i to 480p while other games, such as Gran Turismo 4 and Tourist Trophy are known to support up-scaled 1080i resolution[43] using any of the following standards: composite video[44] (480i), S-Video[45] (480i), RGB[46] (480i/p), VGA[47] (for progressive scan games and PS2 Linux only), YPBPR component video[48] (which display most original PlayStation games in their native 240p mode which most HDTV sets do not support[49]), and D-Terminal[50]. Cables are available for all of these signal types; these cables also output analog stereo audio. Additionally, an RF modulator is available for the system to connect to older TVs[51].
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Digital (S/PDIF) audio may also be output by the console via its TOSLINK connector[52] which outputs 2.0 PCM, 5.1, and 6.1 channel sound in Dolby Digital, Dolby Digital Surround EX, DTS, And DTS-ES formats.
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PlayStation 2 users had the option to play select games over the Internet, using dial-up or a broadband Internet connection. The PlayStation 2 Network Adaptor was required for the original models, while the slim models included networking ports on the console. Instead of having a unified, subscription-based online service like Xbox Live as competitor Microsoft later chose for its Xbox console, online multiplayer functionality on the PlayStation 2 was the responsibility of the game publisher and ran on third-party servers. Many games that supported online play exclusively supported broadband Internet access.
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The PS2 has undergone many revisions[53], some only of internal construction and others involving substantial external changes.
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The PS2 is primarily differentiated between models featuring the original "fat" case design and "slimline" models, which were introduced at the end of 2004. In 2010, the Sony Bravia KDL-22PX300 was made available to consumers. It was a 22" HD-Ready television which incorporated a built-in PlayStation 2.[54][55]
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The PS2 standard color is matte black. Several variations in color were produced in different quantities and regions, including ceramic white, light yellow, metallic blue (aqua), metallic silver, navy (star blue), opaque blue (astral blue), opaque black (midnight black), pearl white, sakura purple, satin gold, satin silver, snow white, super red, transparent blue (ocean blue), and also Limited Edition color Pink, which was distributed in some regions such as Oceania, and parts of Asia.[56][57][58]
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In September 2004, Sony unveiled its third major hardware revision. Available in late October 2004, it was smaller, thinner, and quieter than the original versions and included a built-in Ethernet port (in some markets it also had an integrated modem). Due to its thinner profile, it did not contain the 3.5" expansion bay and therefore did not support the internal hard disk drive. It also lacked an internal power supply until a later revision (excluding the Japan version), similar to the GameCube, and had a modified Multitap expansion. The removal of the expansion bay was criticized as a limitation due to the existence of titles such as Final Fantasy XI, which required the use of the HDD.
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Sony also manufactured a consumer device called the PSX that can be used as a digital video recorder and DVD burner in addition to playing PS2 games. The device was released in Japan on December 13, 2003, and was the first Sony product to include the XrossMediaBar interface. It did not sell well in the Japanese market and was not released anywhere else.[59]
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A class-action lawsuit was filed against Sony Computer Entertainment America Inc. on July 16, 2002, in the Superior Court of California, County of San Mateo. The lawsuit addresses consumer reports of inappropriate "no disc" error (disc read error) messages and other problems associated with playing DVDs and CDs on the PlayStation 2.
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Sony settled its "disc read error" lawsuit by compensating the affected customers with US$25, a free game from a specified list, and the reduced cost repair or replacement (at SCEA's discretion) of the damaged system. This settlement was subject to the courts' approval, and hearings began in the US and Canada on April 28, 2006, and May 11, 2006, respectively.[60]
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PlayStation 2 software is distributed on CD-ROM and DVD-ROM; the two formats are differentiated by the color of their discs' bottoms, with CD-ROMs being blue and DVD-ROMs being silver. The PlayStation 2 offered some particularly high-profile exclusive games. Most main entries in the Grand Theft Auto, Final Fantasy, and Metal Gear Solid series were released exclusively for the console. Several prolific series got their start on the PlayStation 2, including God of War, Ratchet & Clank, Jak and Daxter, Devil May Cry, Kingdom Hearts, and Sly Cooper. Grand Theft Auto: San Andreas was the best-selling game on the console.
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Game releases peaked in 2004, but declined with the release of the PlayStation 3 in 2006. The last new games for the console were Final Fantasy XI: Seekers of Adoulin in Asia, FIFA 13 in North America, and Pro Evolution Soccer 2014 in Europe. As of June 30, 2007, a total of 10,035 software titles had been released worldwide including games released in multiple regions as separate titles.[61]
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Initial reviews in 2000 of the PlayStation 2 acclaimed the console, with reviewers commending its hardware and graphics capabilities, its ability to play DVDs, and the system's backwards compatibility with games and hardware for the original PlayStation. Early points of criticism included the lack of online support at the time, its inclusion of only two controller ports, and the system's price at launch compared to the Dreamcast in 2000.[62][63] PC Magazine in 2001 called the console "outstanding", praising its "noteworthy components" such as the Emotion Engine CPU, 32MB of RAM, support for IEEE 1394 (branded as "i.LINK" by Sony and "FireWire" by Apple), and the console's two USB ports while criticizing its "expensive" games and its support for only two controllers without the multitap accessory.[64]
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Later reviews, especially after the launch of the competing GameCube and Xbox systems, continued to praise the PlayStation 2's large game library and DVD playback, while routinely criticizing the PlayStation 2's lesser graphics performance compared to the newer systems and its rudimentary online service compared to Xbox Live. In 2002, CNET rated the console 7.3 out of 10, calling it a "safe bet" despite not being the "newest or most powerful", noting that the console "yields in-game graphics with more jagged edges". CNET also criticized the DVD playback functionality, claiming that the console's video quality was "passable" and that the playback controls were "rudimentary", recommending users to purchase a remote control. The console's two controller ports and expensiveness of its memory cards were also a point of criticism.[65]
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The slim model of the PlayStation 2 received positive reviews, especially for its incredibly small size and built-in networking. The slim console's requirement for a separate power adapter was often criticized while the top-loading disc drive was often noted as being far less likely to break compared to the tray-loading drive of the original model.[66][67]
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Demand for the PlayStation 2 remained strong throughout much of its lifespan, selling over 1.4 million units in Japan by March 31, 2000. Over 10.6 million units were sold worldwide by March 31, 2001.[68] In 2005, the PlayStation 2 became the fastest game console to reach 100 million units shipped, accomplishing the feat within 5 years and 9 months from its launch; this was surpassed 4 years later when the Nintendo DS reached 100 million shipments in 4 years and 5 months from its launch.[69] By July 2009, the system had sold 138.8 million units worldwide, with 51 million of those units sold in PAL regions.[70]
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Overall, over 155 million PlayStation 2 units were sold worldwide by March 31, 2012, the year Sony officially stopped supplying updated sales numbers of the system.[71]
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The PlayStation 2's DualShock 2 controller is largely identical to the PlayStation's DualShock, with the same basic functionality. However, it includes analog pressure sensitivity on the face, shoulder and D-pad buttons, replacing the digital buttons of the original.[72] (These buttons later became digital again with the release of the DualShock 4.[73]) Like its predecessor, the DualShock 2 controller has force feedback, or "vibration" functionality. It is lighter and includes two more levels of vibration.
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Optional hardware includes additional DualShock or DualShock 2 controllers, a PS2 DVD remote control, an internal or external hard disk drive (HDD), a network adapter, horizontal and vertical stands, PlayStation or PS2 memory cards, the multitap for PlayStation or PS2, a USB motion camera (EyeToy), a USB keyboard and mouse, and a headset.
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The original PS2 multitap (SCPH-10090) cannot be plugged into the newer slim models, as the multitap connects to the memory card slot as well as the controller slot and the memory card slot on the slimline is shallower. New slim-design multitaps (SCPH-70120) were manufactured for these models, however third-party adapters also exist to permit original multitaps to be used.
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Early versions of the PS2 could be networked via an i.LINK port, though this had little game support and was dropped. Some third party manufacturers have created devices that allow disabled people to access the PS2 through ordinary switches, etc.
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Some third-party companies, such as JoyTech, have produced LCD monitor and speaker attachments for the PS2, which attach to the back of the console. These allow users to play games without access to a television as long as there is access to mains electricity or a similar power source. These screens can fold down onto the PS2 in a similar fashion to laptop screens.
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There are many accessories for musical games, such as dance pads for Dance Dance Revolution, In the Groove, and Pump It Up titles and High School Musical 3: Senior Year Dance. Konami microphones for use with the Karaoke Revolution games, dual microphones (sold with and used exclusively for SingStar games), various "guitar" controllers (for the Guitar Freaks series and Guitar Hero series), the drum set controller (sold in a box set (or by itself) with a "guitar" controller and a USB microphone (for use with Rock Band and Guitar Hero series, World Tour and newer), and a taiko drum controller for Taiko: Drum Master.
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Specialized controllers include light guns (GunCon), fishing rod and reel controllers, a Dragon Quest VIII "slime" controller, a Final Fantasy X-2 "Tiny Bee" dual pistol controller, an Onimusha 3 katana controller, and a Resident Evil 4 chainsaw controller.
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Unlike the PlayStation, which requires the use of an official Sony PlayStation Mouse to play mouse-compatible games, the few PS2 games with mouse support work with a standard USB mouse as well as a USB trackball.[74] In addition, some of these games also support the usage of a USB keyboard for text input, game control (in lieu of a DualShock or DualShock 2 gamepad, in tandem with a USB mouse), or both.
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Using homebrew programs, it is possible to play various audio and video file formats on a PS2. Homebrew programs can also be used to play patched backups of original PS2 DVD games on unmodified consoles, and to install retail discs to an installed hard drive on older models. Homebrew emulators of older computer and gaming systems have been developed for the PS2.[75]
|
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|
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Sony released a Linux-based operating system, Linux for PlayStation 2, for the PS2 in a package that also includes a keyboard, mouse, Ethernet adapter and HDD. In Europe and Australia, the PS2 comes with a free Yabasic interpreter on the bundled demo disc. This allows users to create simple programs for the PS2. A port of the NetBSD project and BlackRhino GNU/Linux, an alternative Debian-based distribution, are also available for the PS2.
|
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A successor, the PlayStation 3 was released in Japan and North America in November 2006 and Europe in March 2007.
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1 |
+
A pseudonym (/ˈsjuːdənɪm/) or alias (/ˈeɪliəs/) is a name that a person or group assumes for a particular purpose, which can differ from their first or true name (orthonym).[1] The term is not used when a new name entirely replaces an individual's own.
|
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Pseudonyms include stage names and user names, ring names, pen names, nicknames, aliases, superhero or villain identities and code names, gamer identifications, and regnal names of emperors, popes, and other monarchs. Historically, they have sometimes taken the form of anagrams, Graecisms, and Latinisations, although there are many other methods of choosing a pseudonym.[2]
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Pseudonyms should not be confused with new names that replace old ones and become the individual's full-time name. Pseudonyms are "part-time" names, used only in certain contexts – to provide a more clear-cut separation between one's private and professional lives, to showcase or enhance a particular persona, or to hide an individual's real identity, as with writers' pen names, graffiti artists' tags, resistance fighters' or terrorists' noms de guerre, and computer hackers' handles. Actors, voice-over artists, musicians, and other performers sometimes use stage names, for example, to better channel a relevant energy, gain a greater sense of security and comfort via privacy, more easily avoid troublesome fans/"stalkers", or to mask their ethnic backgrounds.
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In some cases, pseudonyms are adopted because they are part of a cultural or organisational tradition: for example devotional names used by members of some religious institutes, and "cadre names" used by Communist party leaders such as Trotsky and Lenin.
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A pseudonym may also be used for personal reasons: for example, an individual may prefer to be called or known by a name that differs from their given or legal name, but is not ready to take the numerous steps to get their name legally changed; or an individual may simply feel that the context and content of an exchange offer no reason, legal or otherwise, to provide their given or legal name.
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A collective name or collective pseudonym is one shared by two or more persons, for example the co-authors of a work, such as Carolyn Keene, Ellery Queen, Nicolas Bourbaki, or James S. A. Corey.
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The term pseudonym is derived from the Greek ψευδώνυμον (pseudṓnymon), literally "false name", from ψεῦδος (pseûdos), "lie, falsehood"[3] and ὄνομα (ónoma), "name".[4] The term alias is a Latin adverb meaning "at another time, elsewhere".[5]
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A pseudonym is distinct from an allonym, which is the (real) name of another person, assumed by the author of a work of art.[6] This may occur when someone is ghostwriting a book or play, or in parody, or when using a "front" name, such as by screenwriters blacklisted in Hollywood in the 1950s and 1960s. See also pseudepigraph, for falsely attributed authorship.
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Sometimes people change their name in such a manner that the new name becomes permanent and is used by all who know the person. This is not an alias or pseudonym, but in fact a new name. In many countries, including common law countries, a name change can be ratified by a court and become a person's new legal name.
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For example, in the 1960s, black civil rights campaigner Malcolm Little changed his surname to "X", to represent his unknown African ancestral name that had been lost when his ancestors were brought to North America as slaves. He then changed his name again to Malik El-Shabazz when he converted to Islam.[citation needed] Likewise some Jews adopted Hebrew family names upon immigrating to Israel, dropping surnames that had been in their families for generations. The politician David Ben-Gurion, for example, was born David Grün in Poland. He adopted his Hebrew name in 1910, when he published his first article in a Zionist journal in Jerusalem.[7] Many transgender people also choose to adopt a new name, typically around the time of their social transitioning, to match their desired gender better than their birth name.[according to whom?]
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Businesspersons of ethnic minorities in some parts of the world are sometimes advised by an employer to use a pseudonym that is common or acceptable in that area when conducting business, to overcome racial or religious bias.[8]
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Criminals may use aliases, fictitious business names, and dummy corporations (corporate shells) to hide their identity, or to impersonate other persons or entities in order to commit fraud. Aliases and fictitious business names used for dummy corporations may become so complex that, in the words of the Washington Post, "getting to the truth requires a walk down a bizarre labyrinth" and multiple government agencies may become involved to uncover the truth.[9]
|
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A pen name, or "nom de plume" (French for "pen name"), is a pseudonym (sometimes a particular form of the real name) adopted by an author (or on the author's behalf by their publishers).
|
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|
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+
Some female authors used male pen names, in particular in the 19th century, when writing was a male-dominated profession. The Brontë sisters used pen names for their early work, so as not to reveal their gender (see below) and so that local residents would not know that the books related to people of the neighbourhood. The Brontës used their neighbours as inspiration for characters in many of their books. Anne Brontë's The Tenant of Wildfell Hall (1848) was published under the name Acton Bell, while Charlotte Brontë used the name Currer Bell for Jane Eyre (1847) and Shirley (1849), and Emily Brontë adopted Ellis Bell as cover for Wuthering Heights (1847). Other examples from the nineteenth-century are the novelist Mary Ann Evans (George Eliot) and the French writer Amandine Aurore Lucile Dupin (George Sand). Pseudonyms may also be used due to cultural or organization or political prejudices.
|
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|
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On the other hand, some 20th and 21st-century male romance novelists have used female pen names.[10] A few examples are Brindle Chase, Peter O'Donnell (as Madeline Brent), Christopher Wood (as Penny Sutton and Rosie Dixon), and Hugh C. Rae (as Jessica Sterling).[10]
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A pen name may be used if a writer's real name is likely to be confused with the name of another writer or notable individual, or if the real name is deemed unsuitable.
|
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Authors who write both fiction and non-fiction, or in different genres, may use different pen names to avoid confusing their readers. For example, the romance writer Nora Roberts writes mystery novels under the name J.D. Robb.
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In some cases, an author may become better known by his pen name than his real name. A famous example is Samuel Clemens, writing as Mark Twain. The British mathematician Charles Dodgson wrote fantasy novels as Lewis Carroll and mathematical treatises under his own name.
|
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|
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+
Some authors, such as Harold Robbins, use several literary pseudonyms.[11]
|
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|
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+
Some pen names have been used for long periods, even decades, without the author's true identity being discovered, as with Elena Ferrante and Torsten Krol.
|
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|
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Some pen names are not strictly pseudonyms, as they are simply variants of the authors' actual names. The authors C.L. Moore and S.E. Hinton were female authors who used the initialled forms of their full names, Moore being Catherine Lucille Moore, writing in the 1930s male-dominated science fiction genre, and Hinton, (author of The Outsiders) Susan Eloise Hinton. Star Trek writer D.C. Fontana (Dorothy Catherine) wrote using her own abbreviated name and under the pen names Michael Richards and J. Michael Bingham. Author V.C. Andrews intended to publish under her given name, Virginia Andrews, but was told that, due to a production error, her first novel was being released under the name of "V.C. Andrews"; later she learned that her publisher had in fact done this deliberately. Joanne Kathleen Rowling[12] published the Harry Potter series as J.K. Rowling. Rowling also published the Cormoran Strike series, a series of detective novels including The Cuckoo's Calling under the pseudonym "Robert Galbraith".
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Winston Churchill wrote as Winston S. Churchill (from his full surname "Spencer-Churchill" which he did not otherwise use) in an attempt to avoid confusion with an American novelist of the same name. The attempt was not wholly successful – the two are still sometimes confused by booksellers.[13][14]
|
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A pen name may be used specifically to hide the identity of the author, as with exposé books about espionage or crime, or explicit erotic fiction. Some prolific authors adopt a pseudonym to disguise the extent of their published output, e. g. Stephen King writing as Richard Bachman. Co-authors may choose to publish under a collective pseudonym, e. g., P. J. Tracy and Perri O'Shaughnessy. Frederic Dannay and Manfred Lee used the name Ellery Queen as a pen name for their collaborative works and as the name of their main character.[citation needed] Asa Earl Carter, a Southern white segregationist affiliated with the KKK, wrote Western books under a fictional Cherokee persona to imply legitimacy and conceal his history.[15]
|
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"Why do authors choose pseudonyms? It is rarely because they actually hope to stay anonymous forever," mused writer and columnist Russell Smith in his review of the Canadian novel Into That Fire by the pseudonymous M.J. Cates.[16]
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A famous case in French literature was Romain Gary. Already a well-known writer, he started publishing books as Émile Ajar to test whether his new books would be well received on their own merits, without the aid of his established reputation. They were: Émile Ajar, like Romain Gary before him, was awarded the prestigious Prix Goncourt by a jury unaware that they were the same person. Similarly, TV actor Ronnie Barker submitted comedy material under the name Gerald Wiley.
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A collective pseudonym may represent an entire publishing house, or any contributor to a long-running series, especially with juvenile literature. Examples include Watty Piper, Victor Appleton, Erin Hunter, and Kamiru M. Xhan.
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Another use of a pseudonym in literature is to present a story as being written by the fictional characters in the story. The series of novels known as A Series Of Unfortunate Events are written by Daniel Handler under the pen name of Lemony Snicket, a character in the series. This applies also to some of the several 18th-century English and American writers who used the name Fidelia.
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An anonymity pseudonym or multiple-use name is a name used by many different people to protect anonymity.[17] It is a strategy that has been adopted by many unconnected radical groups and by cultural groups, where the construct of personal identity has been criticised. This has led to the idea of the "open pop star".
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Pseudonyms and acronyms are often employed in medical research to protect subjects' identities through a process known as de-identification.
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Nicolaus Copernicus put forward his theory of heliocentrism in the manuscript Commentariolus anonymously.[18] Sophie Germain and William Sealy Gosset used pseudonyms to publish their work in the field of mathematics.[19][20] Satoshi Nakamoto is a pseudonym of a still unknown author or authors' group behind a white paper about bitcoin.[21][22][23][24] In the field of physics, one case of usage of pseudonyms is denounced.[25] Ignazio Ciufolini is accused of publishing two papers on the scientific preprint archive arXiv.org under pseudonyms, each criticizing one of the rivals to LAGEOS, what is argued to be a form of ventriloquism.[26] Such conduct is a violation of arXiv terms of use.[27][28][26]
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In Ancien Régime France, a nom de guerre ("war name") would be adopted by each new recruit (or assigned to them by the captain of their company) as they enlisted in the French army. These pseudonyms had an official character and were the predecessor of identification numbers: soldiers were identified by their first names, their family names, and their noms de guerre (e. g. Jean Amarault dit Lafidélité). These pseudonyms were usually related to the soldier's place of origin (e. g. Jean Deslandes dit Champigny, for a soldier coming from a town named Champigny), or to a particular physical or personal trait (e. g. Antoine Bonnet dit Prettaboire, for a soldier prêt à boire, ready to drink). In 1716, a nom de guerre was mandatory for every soldier; officers did not adopt noms de guerre as they considered them derogatory. In daily life, these aliases could replace the real family name.[29]
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Noms de guerre were adopted for security reasons by members of the World War II French resistance and Polish resistance. Such pseudonyms are often adopted by military special-forces soldiers, such as members of the SAS and similar units of resistance fighters, terrorists, and guerrillas. This practice hides their identities and may protect their families from reprisals; it may also be a form of dissociation from domestic life. Some well-known men who adopted noms de guerre include Carlos, for Ilich Ramírez Sánchez; Willy Brandt, Chancellor of West Germany; and Subcomandante Marcos, spokesman of the Zapatista Army of National Liberation (EZLN).[citation needed] During Lehi's underground fight against the British in Mandatory Palestine, the organization's commander Yitzchak Shamir (later Prime Minister of Israel) adopted the nom de guerre "Michael", in honour of Ireland's Michael Collins.
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Revolutionaries and resistance leaders, such as Lenin, Trotsky, Golda Meir, Philippe Leclerc de Hauteclocque, and Josip Broz Tito, often adopted their noms de guerre as their proper names after the struggle. George Grivas, the Greek-Cypriot EOKA militant, adopted the nom de guerre Digenis (Διγενής). In the French Foreign Legion, recruits can adopt a pseudonym to break with their past lives. Mercenaries have long used "noms de guerre", sometimes even multiple identities, depending on country, conflict and circumstance.[citation needed] Some of the most familiar noms de guerre today are the kunya used by Islamic mujahideen. These take the form of a teknonym, either literal or figurative.
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Individuals using a computer online may adopt or be required to use a form of pseudonym known as a "handle" (a term deriving from CB slang), "user name", "login name", "avatar", or, sometimes, "screen name", "gamertag" "IGN (In Game (Nick)Name)" or "nickname". On the Internet, pseudonymous remailers use cryptography that achieves persistent pseudonymity, so that two-way communication can be achieved, and reputations can be established, without linking physical identities to their respective pseudonyms. Aliasing is the use of multiple names for the same data location.
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More sophisticated cryptographic systems, such as anonymous digital credentials, enable users to communicate pseudonymously (i. e., by identifying themselves by means of pseudonyms). In well-defined abuse cases, a designated authority may be able to revoke the pseudonyms and reveal the individuals' real identity.[citation needed]
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Use of pseudonyms is common among professional eSports players, despite the fact that many professional games are played on LAN.[30]
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People seeking privacy often use pseudonyms to make appointments and reservations.[31] Those writing to advice columns in newspapers and magazines may use pseudonyms.[32] Steve Wozniak used a pseudonym when attending the University of California, Berkeley after co-founding Apple Computer, because "I knew I wouldn't have time enough to be an A+ student."[33]
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When used by an actor, musician, radio disc jockey, model, or other performer or "show business" personality a pseudonym is called a stage name, or, occasionally, a professional name, or screen name.
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Members of a marginalized ethnic or religious group have often adopted stage names, typically changing their surname or entire name to mask their original background.
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Stage names are also used to create a more marketable name, as in the case of Creighton Tull Chaney, who adopted the pseudonym Lon Chaney, Jr., a reference to his famous father Lon Chaney, Sr.
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Chris Curtis of Deep Purple fame was christened as Christopher Crummey ("crumby" is UK slang for poor quality). In this and similar cases a stage name is adopted simply to avoid an unfortunate pun.
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Pseudonyms are also used to comply with the rules of performing arts guilds (Screen Actors Guild (SAG), Writers Guild of America, East (WGA), AFTRA, etc.), which do not allow performers to use an existing name, in order to avoid confusion. For example, these rules required film and television actor Michael Fox to add a middle initial and become Michael J. Fox, to avoid being confused with another actor named Michael Fox. This was also true of author and actress Fannie Flagg, who chose this pseudonym; her real name, Patricia Neal, being the name of another well-known actress; and British actor Stewart Granger, whose real name was James Stewart. The film-making team of Joel and Ethan Coen, for instance, share credit for editing under the alias Roderick Jaynes.[34]
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Some stage names are used to conceal a person's identity, such as the pseudonym Alan Smithee, which was used by directors in the Directors Guild of America (DGA) to remove their name from a film they feel was edited or modified beyond their artistic satisfaction. In theatre, the pseudonyms George or Georgina Spelvin, and Walter Plinge are used to hide the identity of a performer, usually when he or she is "doubling" (playing more than one role in the same play).
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David Agnew was a name used by the BBC to conceal the identity of a scriptwriter, such as for the Doctor Who serial City of Death, which had three writers, including Douglas Adams, who was at the time of writing the show's Script Editor.[35] In another Doctor Who serial, The Brain of Morbius, writer Terrance Dicks demanded the removal of his name from the credits saying it could go out under a "bland pseudonym".[citation needed][36] This ended up as Robin Bland.[36][37]
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Musicians and singers can use pseudonyms to allow artists to collaborate with artists on other labels while avoiding the need to gain permission from their own labels, such as the artist Jerry Samuels, who made songs under Napoleon XIV. Rock singer-guitarist George Harrison, for example, played guitar on Cream's song "Badge" using a pseudonym.[38] In classical music, some record companies issued recordings under a nom de disque in the 1950s and 1960s to avoid paying royalties. A number of popular budget LPs of piano music were released under the pseudonym Paul Procopolis.[citation needed] Another example is that Paul McCartney used his fictional name "Bernerd Webb" for Peter and Gordon's song Woman.[39]
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Pseudonyms are used as stage names in heavy metal bands, such as Tracii Guns in LA Guns, Axl Rose and Slash in Guns N' Roses, Mick Mars in Mötley Crüe, Dimebag Darrell in Pantera, or C.C. Deville in Poison. Some such names have additional meanings, like that of Brian Hugh Warner, more commonly known as Marilyn Manson: Marilyn coming from Marilyn Monroe and Manson from convicted serial killer Charles Manson. Jacoby Shaddix of Papa Roach went under the name "Coby Dick" during the Infest era. He changed back to his birth name when lovehatetragedy was released.
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David Johansen, front man for the hard rock band New York Dolls, recorded and performed pop and lounge music under the pseudonym Buster Poindexter in the late 1980s and early 1990s. The music video for Poindexter's debt single, Hot Hot Hot, opens with a monologue from Johansen where he notes his time with the New York Dolls and explains his desire to create more sophisticated music.
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Ross Bagdasarian, Sr., creator of Alvin and the Chipmunks, wrote original songs, arranged, and produced the records under his real name, but performed on them as David Seville. He also wrote songs as Skipper Adams. Danish pop pianist Bent Fabric, whose full name is Bent Fabricius-Bjerre, wrote his biggest instrumental hit "Alley Cat" as Frank Bjorn.
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For a time, the musician Prince used an unpronounceable "Love Symbol" as a pseudonym ("Prince" is his actual first name rather than a stage name). He wrote the song "Sugar Walls" for Sheena Easton as "Alexander Nevermind" and "Manic Monday" for The Bangles as "Christopher Tracy". (He also produced albums early in his career as "Jamie Starr").
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Many Italian-American singers have used stage names, as their birth names were difficult to pronounce or considered too ethnic for American tastes. Singers changing their names included Dean Martin (born Dino Paul Crocetti), Connie Francis (born Concetta Franconero), Frankie Valli (born Francesco Castelluccio), Tony Bennett (born Anthony Benedetto), and Lady Gaga (born Stefani Germanotta)
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In 2009, the British rock band Feeder briefly changed its name to Renegades so it could play a whole show featuring a set list in which 95 per cent of the songs played were from their forthcoming new album of the same name, with none of their singles included. Front man Grant Nicholas felt that if they played as Feeder, there would be uproar over him not playing any of the singles, so used the pseudonym as a hint. A series of small shows were played in 2010, at 250- to 1,000-capacity venues with the plan not to say who the band really are and just announce the shows as if they were a new band.
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In many cases, hip-hop and rap artists prefer to use pseudonyms that represents some variation of their name, personality, or interests. Examples include Iggy Azalea (her stage name is a combination of her dog's name, Iggy, and her home street in Mullumbimby, Azalea Street), Ol' Dirty Bastard (known under at least six aliases), Diddy (previously known at various times as Puffy, P. Diddy, and Puff Daddy), Ludacris, Flo Rida (whose stage name is a tribute to his home state, Florida), British-Jamaican hip-hop artist Stefflon Don (real name Stephanie Victoria Allen), LL Cool J, and Chingy. Black metal artists also adopt pseudonyms, usually symbolizing dark values, such as Nocturno Culto, Gaahl, Abbath, and Silenoz. In punk and hardcore punk, singers and band members often replace real names with tougher-sounding stage names such as Sid Vicious (real name John Simon Ritchie) of the late 1970s band Sex Pistols and "Rat" of the early 1980s band The Varukers and the 2000s re-formation of Discharge. The punk rock band The Ramones had every member take the last name of Ramone.[citation needed]
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Henry John Deutschendorf Jr., an American singer-songwriter, used the stage name John Denver. The Australian country musician born Robert Lane changed his name to Tex Morton. Reginald Kenneth Dwight legally changed his name in 1972 to Elton John.
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A pseudonym (/ˈsjuːdənɪm/) or alias (/ˈeɪliəs/) is a name that a person or group assumes for a particular purpose, which can differ from their first or true name (orthonym).[1] The term is not used when a new name entirely replaces an individual's own.
|
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Pseudonyms include stage names and user names, ring names, pen names, nicknames, aliases, superhero or villain identities and code names, gamer identifications, and regnal names of emperors, popes, and other monarchs. Historically, they have sometimes taken the form of anagrams, Graecisms, and Latinisations, although there are many other methods of choosing a pseudonym.[2]
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Pseudonyms should not be confused with new names that replace old ones and become the individual's full-time name. Pseudonyms are "part-time" names, used only in certain contexts – to provide a more clear-cut separation between one's private and professional lives, to showcase or enhance a particular persona, or to hide an individual's real identity, as with writers' pen names, graffiti artists' tags, resistance fighters' or terrorists' noms de guerre, and computer hackers' handles. Actors, voice-over artists, musicians, and other performers sometimes use stage names, for example, to better channel a relevant energy, gain a greater sense of security and comfort via privacy, more easily avoid troublesome fans/"stalkers", or to mask their ethnic backgrounds.
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In some cases, pseudonyms are adopted because they are part of a cultural or organisational tradition: for example devotional names used by members of some religious institutes, and "cadre names" used by Communist party leaders such as Trotsky and Lenin.
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A pseudonym may also be used for personal reasons: for example, an individual may prefer to be called or known by a name that differs from their given or legal name, but is not ready to take the numerous steps to get their name legally changed; or an individual may simply feel that the context and content of an exchange offer no reason, legal or otherwise, to provide their given or legal name.
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A collective name or collective pseudonym is one shared by two or more persons, for example the co-authors of a work, such as Carolyn Keene, Ellery Queen, Nicolas Bourbaki, or James S. A. Corey.
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The term pseudonym is derived from the Greek ψευδώνυμον (pseudṓnymon), literally "false name", from ψεῦδος (pseûdos), "lie, falsehood"[3] and ὄνομα (ónoma), "name".[4] The term alias is a Latin adverb meaning "at another time, elsewhere".[5]
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A pseudonym is distinct from an allonym, which is the (real) name of another person, assumed by the author of a work of art.[6] This may occur when someone is ghostwriting a book or play, or in parody, or when using a "front" name, such as by screenwriters blacklisted in Hollywood in the 1950s and 1960s. See also pseudepigraph, for falsely attributed authorship.
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Sometimes people change their name in such a manner that the new name becomes permanent and is used by all who know the person. This is not an alias or pseudonym, but in fact a new name. In many countries, including common law countries, a name change can be ratified by a court and become a person's new legal name.
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For example, in the 1960s, black civil rights campaigner Malcolm Little changed his surname to "X", to represent his unknown African ancestral name that had been lost when his ancestors were brought to North America as slaves. He then changed his name again to Malik El-Shabazz when he converted to Islam.[citation needed] Likewise some Jews adopted Hebrew family names upon immigrating to Israel, dropping surnames that had been in their families for generations. The politician David Ben-Gurion, for example, was born David Grün in Poland. He adopted his Hebrew name in 1910, when he published his first article in a Zionist journal in Jerusalem.[7] Many transgender people also choose to adopt a new name, typically around the time of their social transitioning, to match their desired gender better than their birth name.[according to whom?]
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Businesspersons of ethnic minorities in some parts of the world are sometimes advised by an employer to use a pseudonym that is common or acceptable in that area when conducting business, to overcome racial or religious bias.[8]
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Criminals may use aliases, fictitious business names, and dummy corporations (corporate shells) to hide their identity, or to impersonate other persons or entities in order to commit fraud. Aliases and fictitious business names used for dummy corporations may become so complex that, in the words of the Washington Post, "getting to the truth requires a walk down a bizarre labyrinth" and multiple government agencies may become involved to uncover the truth.[9]
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A pen name, or "nom de plume" (French for "pen name"), is a pseudonym (sometimes a particular form of the real name) adopted by an author (or on the author's behalf by their publishers).
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Some female authors used male pen names, in particular in the 19th century, when writing was a male-dominated profession. The Brontë sisters used pen names for their early work, so as not to reveal their gender (see below) and so that local residents would not know that the books related to people of the neighbourhood. The Brontës used their neighbours as inspiration for characters in many of their books. Anne Brontë's The Tenant of Wildfell Hall (1848) was published under the name Acton Bell, while Charlotte Brontë used the name Currer Bell for Jane Eyre (1847) and Shirley (1849), and Emily Brontë adopted Ellis Bell as cover for Wuthering Heights (1847). Other examples from the nineteenth-century are the novelist Mary Ann Evans (George Eliot) and the French writer Amandine Aurore Lucile Dupin (George Sand). Pseudonyms may also be used due to cultural or organization or political prejudices.
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On the other hand, some 20th and 21st-century male romance novelists have used female pen names.[10] A few examples are Brindle Chase, Peter O'Donnell (as Madeline Brent), Christopher Wood (as Penny Sutton and Rosie Dixon), and Hugh C. Rae (as Jessica Sterling).[10]
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A pen name may be used if a writer's real name is likely to be confused with the name of another writer or notable individual, or if the real name is deemed unsuitable.
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Authors who write both fiction and non-fiction, or in different genres, may use different pen names to avoid confusing their readers. For example, the romance writer Nora Roberts writes mystery novels under the name J.D. Robb.
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In some cases, an author may become better known by his pen name than his real name. A famous example is Samuel Clemens, writing as Mark Twain. The British mathematician Charles Dodgson wrote fantasy novels as Lewis Carroll and mathematical treatises under his own name.
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Some authors, such as Harold Robbins, use several literary pseudonyms.[11]
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Some pen names have been used for long periods, even decades, without the author's true identity being discovered, as with Elena Ferrante and Torsten Krol.
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Some pen names are not strictly pseudonyms, as they are simply variants of the authors' actual names. The authors C.L. Moore and S.E. Hinton were female authors who used the initialled forms of their full names, Moore being Catherine Lucille Moore, writing in the 1930s male-dominated science fiction genre, and Hinton, (author of The Outsiders) Susan Eloise Hinton. Star Trek writer D.C. Fontana (Dorothy Catherine) wrote using her own abbreviated name and under the pen names Michael Richards and J. Michael Bingham. Author V.C. Andrews intended to publish under her given name, Virginia Andrews, but was told that, due to a production error, her first novel was being released under the name of "V.C. Andrews"; later she learned that her publisher had in fact done this deliberately. Joanne Kathleen Rowling[12] published the Harry Potter series as J.K. Rowling. Rowling also published the Cormoran Strike series, a series of detective novels including The Cuckoo's Calling under the pseudonym "Robert Galbraith".
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Winston Churchill wrote as Winston S. Churchill (from his full surname "Spencer-Churchill" which he did not otherwise use) in an attempt to avoid confusion with an American novelist of the same name. The attempt was not wholly successful – the two are still sometimes confused by booksellers.[13][14]
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A pen name may be used specifically to hide the identity of the author, as with exposé books about espionage or crime, or explicit erotic fiction. Some prolific authors adopt a pseudonym to disguise the extent of their published output, e. g. Stephen King writing as Richard Bachman. Co-authors may choose to publish under a collective pseudonym, e. g., P. J. Tracy and Perri O'Shaughnessy. Frederic Dannay and Manfred Lee used the name Ellery Queen as a pen name for their collaborative works and as the name of their main character.[citation needed] Asa Earl Carter, a Southern white segregationist affiliated with the KKK, wrote Western books under a fictional Cherokee persona to imply legitimacy and conceal his history.[15]
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"Why do authors choose pseudonyms? It is rarely because they actually hope to stay anonymous forever," mused writer and columnist Russell Smith in his review of the Canadian novel Into That Fire by the pseudonymous M.J. Cates.[16]
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A famous case in French literature was Romain Gary. Already a well-known writer, he started publishing books as Émile Ajar to test whether his new books would be well received on their own merits, without the aid of his established reputation. They were: Émile Ajar, like Romain Gary before him, was awarded the prestigious Prix Goncourt by a jury unaware that they were the same person. Similarly, TV actor Ronnie Barker submitted comedy material under the name Gerald Wiley.
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A collective pseudonym may represent an entire publishing house, or any contributor to a long-running series, especially with juvenile literature. Examples include Watty Piper, Victor Appleton, Erin Hunter, and Kamiru M. Xhan.
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Another use of a pseudonym in literature is to present a story as being written by the fictional characters in the story. The series of novels known as A Series Of Unfortunate Events are written by Daniel Handler under the pen name of Lemony Snicket, a character in the series. This applies also to some of the several 18th-century English and American writers who used the name Fidelia.
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An anonymity pseudonym or multiple-use name is a name used by many different people to protect anonymity.[17] It is a strategy that has been adopted by many unconnected radical groups and by cultural groups, where the construct of personal identity has been criticised. This has led to the idea of the "open pop star".
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Pseudonyms and acronyms are often employed in medical research to protect subjects' identities through a process known as de-identification.
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Nicolaus Copernicus put forward his theory of heliocentrism in the manuscript Commentariolus anonymously.[18] Sophie Germain and William Sealy Gosset used pseudonyms to publish their work in the field of mathematics.[19][20] Satoshi Nakamoto is a pseudonym of a still unknown author or authors' group behind a white paper about bitcoin.[21][22][23][24] In the field of physics, one case of usage of pseudonyms is denounced.[25] Ignazio Ciufolini is accused of publishing two papers on the scientific preprint archive arXiv.org under pseudonyms, each criticizing one of the rivals to LAGEOS, what is argued to be a form of ventriloquism.[26] Such conduct is a violation of arXiv terms of use.[27][28][26]
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In Ancien Régime France, a nom de guerre ("war name") would be adopted by each new recruit (or assigned to them by the captain of their company) as they enlisted in the French army. These pseudonyms had an official character and were the predecessor of identification numbers: soldiers were identified by their first names, their family names, and their noms de guerre (e. g. Jean Amarault dit Lafidélité). These pseudonyms were usually related to the soldier's place of origin (e. g. Jean Deslandes dit Champigny, for a soldier coming from a town named Champigny), or to a particular physical or personal trait (e. g. Antoine Bonnet dit Prettaboire, for a soldier prêt à boire, ready to drink). In 1716, a nom de guerre was mandatory for every soldier; officers did not adopt noms de guerre as they considered them derogatory. In daily life, these aliases could replace the real family name.[29]
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Noms de guerre were adopted for security reasons by members of the World War II French resistance and Polish resistance. Such pseudonyms are often adopted by military special-forces soldiers, such as members of the SAS and similar units of resistance fighters, terrorists, and guerrillas. This practice hides their identities and may protect their families from reprisals; it may also be a form of dissociation from domestic life. Some well-known men who adopted noms de guerre include Carlos, for Ilich Ramírez Sánchez; Willy Brandt, Chancellor of West Germany; and Subcomandante Marcos, spokesman of the Zapatista Army of National Liberation (EZLN).[citation needed] During Lehi's underground fight against the British in Mandatory Palestine, the organization's commander Yitzchak Shamir (later Prime Minister of Israel) adopted the nom de guerre "Michael", in honour of Ireland's Michael Collins.
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Revolutionaries and resistance leaders, such as Lenin, Trotsky, Golda Meir, Philippe Leclerc de Hauteclocque, and Josip Broz Tito, often adopted their noms de guerre as their proper names after the struggle. George Grivas, the Greek-Cypriot EOKA militant, adopted the nom de guerre Digenis (Διγενής). In the French Foreign Legion, recruits can adopt a pseudonym to break with their past lives. Mercenaries have long used "noms de guerre", sometimes even multiple identities, depending on country, conflict and circumstance.[citation needed] Some of the most familiar noms de guerre today are the kunya used by Islamic mujahideen. These take the form of a teknonym, either literal or figurative.
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Individuals using a computer online may adopt or be required to use a form of pseudonym known as a "handle" (a term deriving from CB slang), "user name", "login name", "avatar", or, sometimes, "screen name", "gamertag" "IGN (In Game (Nick)Name)" or "nickname". On the Internet, pseudonymous remailers use cryptography that achieves persistent pseudonymity, so that two-way communication can be achieved, and reputations can be established, without linking physical identities to their respective pseudonyms. Aliasing is the use of multiple names for the same data location.
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More sophisticated cryptographic systems, such as anonymous digital credentials, enable users to communicate pseudonymously (i. e., by identifying themselves by means of pseudonyms). In well-defined abuse cases, a designated authority may be able to revoke the pseudonyms and reveal the individuals' real identity.[citation needed]
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Use of pseudonyms is common among professional eSports players, despite the fact that many professional games are played on LAN.[30]
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People seeking privacy often use pseudonyms to make appointments and reservations.[31] Those writing to advice columns in newspapers and magazines may use pseudonyms.[32] Steve Wozniak used a pseudonym when attending the University of California, Berkeley after co-founding Apple Computer, because "I knew I wouldn't have time enough to be an A+ student."[33]
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When used by an actor, musician, radio disc jockey, model, or other performer or "show business" personality a pseudonym is called a stage name, or, occasionally, a professional name, or screen name.
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Members of a marginalized ethnic or religious group have often adopted stage names, typically changing their surname or entire name to mask their original background.
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Stage names are also used to create a more marketable name, as in the case of Creighton Tull Chaney, who adopted the pseudonym Lon Chaney, Jr., a reference to his famous father Lon Chaney, Sr.
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Chris Curtis of Deep Purple fame was christened as Christopher Crummey ("crumby" is UK slang for poor quality). In this and similar cases a stage name is adopted simply to avoid an unfortunate pun.
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Pseudonyms are also used to comply with the rules of performing arts guilds (Screen Actors Guild (SAG), Writers Guild of America, East (WGA), AFTRA, etc.), which do not allow performers to use an existing name, in order to avoid confusion. For example, these rules required film and television actor Michael Fox to add a middle initial and become Michael J. Fox, to avoid being confused with another actor named Michael Fox. This was also true of author and actress Fannie Flagg, who chose this pseudonym; her real name, Patricia Neal, being the name of another well-known actress; and British actor Stewart Granger, whose real name was James Stewart. The film-making team of Joel and Ethan Coen, for instance, share credit for editing under the alias Roderick Jaynes.[34]
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Some stage names are used to conceal a person's identity, such as the pseudonym Alan Smithee, which was used by directors in the Directors Guild of America (DGA) to remove their name from a film they feel was edited or modified beyond their artistic satisfaction. In theatre, the pseudonyms George or Georgina Spelvin, and Walter Plinge are used to hide the identity of a performer, usually when he or she is "doubling" (playing more than one role in the same play).
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David Agnew was a name used by the BBC to conceal the identity of a scriptwriter, such as for the Doctor Who serial City of Death, which had three writers, including Douglas Adams, who was at the time of writing the show's Script Editor.[35] In another Doctor Who serial, The Brain of Morbius, writer Terrance Dicks demanded the removal of his name from the credits saying it could go out under a "bland pseudonym".[citation needed][36] This ended up as Robin Bland.[36][37]
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Musicians and singers can use pseudonyms to allow artists to collaborate with artists on other labels while avoiding the need to gain permission from their own labels, such as the artist Jerry Samuels, who made songs under Napoleon XIV. Rock singer-guitarist George Harrison, for example, played guitar on Cream's song "Badge" using a pseudonym.[38] In classical music, some record companies issued recordings under a nom de disque in the 1950s and 1960s to avoid paying royalties. A number of popular budget LPs of piano music were released under the pseudonym Paul Procopolis.[citation needed] Another example is that Paul McCartney used his fictional name "Bernerd Webb" for Peter and Gordon's song Woman.[39]
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Pseudonyms are used as stage names in heavy metal bands, such as Tracii Guns in LA Guns, Axl Rose and Slash in Guns N' Roses, Mick Mars in Mötley Crüe, Dimebag Darrell in Pantera, or C.C. Deville in Poison. Some such names have additional meanings, like that of Brian Hugh Warner, more commonly known as Marilyn Manson: Marilyn coming from Marilyn Monroe and Manson from convicted serial killer Charles Manson. Jacoby Shaddix of Papa Roach went under the name "Coby Dick" during the Infest era. He changed back to his birth name when lovehatetragedy was released.
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David Johansen, front man for the hard rock band New York Dolls, recorded and performed pop and lounge music under the pseudonym Buster Poindexter in the late 1980s and early 1990s. The music video for Poindexter's debt single, Hot Hot Hot, opens with a monologue from Johansen where he notes his time with the New York Dolls and explains his desire to create more sophisticated music.
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Ross Bagdasarian, Sr., creator of Alvin and the Chipmunks, wrote original songs, arranged, and produced the records under his real name, but performed on them as David Seville. He also wrote songs as Skipper Adams. Danish pop pianist Bent Fabric, whose full name is Bent Fabricius-Bjerre, wrote his biggest instrumental hit "Alley Cat" as Frank Bjorn.
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For a time, the musician Prince used an unpronounceable "Love Symbol" as a pseudonym ("Prince" is his actual first name rather than a stage name). He wrote the song "Sugar Walls" for Sheena Easton as "Alexander Nevermind" and "Manic Monday" for The Bangles as "Christopher Tracy". (He also produced albums early in his career as "Jamie Starr").
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Many Italian-American singers have used stage names, as their birth names were difficult to pronounce or considered too ethnic for American tastes. Singers changing their names included Dean Martin (born Dino Paul Crocetti), Connie Francis (born Concetta Franconero), Frankie Valli (born Francesco Castelluccio), Tony Bennett (born Anthony Benedetto), and Lady Gaga (born Stefani Germanotta)
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In 2009, the British rock band Feeder briefly changed its name to Renegades so it could play a whole show featuring a set list in which 95 per cent of the songs played were from their forthcoming new album of the same name, with none of their singles included. Front man Grant Nicholas felt that if they played as Feeder, there would be uproar over him not playing any of the singles, so used the pseudonym as a hint. A series of small shows were played in 2010, at 250- to 1,000-capacity venues with the plan not to say who the band really are and just announce the shows as if they were a new band.
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In many cases, hip-hop and rap artists prefer to use pseudonyms that represents some variation of their name, personality, or interests. Examples include Iggy Azalea (her stage name is a combination of her dog's name, Iggy, and her home street in Mullumbimby, Azalea Street), Ol' Dirty Bastard (known under at least six aliases), Diddy (previously known at various times as Puffy, P. Diddy, and Puff Daddy), Ludacris, Flo Rida (whose stage name is a tribute to his home state, Florida), British-Jamaican hip-hop artist Stefflon Don (real name Stephanie Victoria Allen), LL Cool J, and Chingy. Black metal artists also adopt pseudonyms, usually symbolizing dark values, such as Nocturno Culto, Gaahl, Abbath, and Silenoz. In punk and hardcore punk, singers and band members often replace real names with tougher-sounding stage names such as Sid Vicious (real name John Simon Ritchie) of the late 1970s band Sex Pistols and "Rat" of the early 1980s band The Varukers and the 2000s re-formation of Discharge. The punk rock band The Ramones had every member take the last name of Ramone.[citation needed]
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Henry John Deutschendorf Jr., an American singer-songwriter, used the stage name John Denver. The Australian country musician born Robert Lane changed his name to Tex Morton. Reginald Kenneth Dwight legally changed his name in 1972 to Elton John.
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A psychologist is a person who studies normal and abnormal mental states, perceptual, cognitive, emotional, and social processes and behavior by experimenting with, and observing, interpreting, and recording how individuals relate to one another and to their environments.[1]
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Applied psychology applies theory to solve problems in human and animal behavior. Clinical psychology is a field of applied psychology that focus on therapeutic methods. Other applied fields include counseling psychology and school psychology.[2] Licensing and regulations can vary by state and profession.[3]
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In the United States and Canada, full membership in the American Psychological Association requires doctoral training (except in some Canadian provinces, such as Alberta, where a master's degree is sufficient).[a] The minimal requirement for full membership can be waived in circumstances where there is evidence that significant contribution or performance in the field of psychology has been made. Associate membership requires at least two years of postgraduate studies in psychology or an approved related discipline.[4]
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There are a number of U.S. schools offering accredited programs in clinical psychology resulting in a master's degree. Such programs can range from forty-eight to eighty-four units, most often taking two to three years to complete after the undergraduate degree. Training usually emphasizes theory and treatment over research, quite often with a focus on school, or couples and family counseling. Similar to doctoral programs, master's level students usually must fulfill time in a clinical practicum under supervision; some programs also require a minimum amount of personal psychotherapy.[5] While many graduates from master's level training go on to doctoral psychology programs, a large number also go directly into practice—often as a licensed professional counselor (LPC), marriage and family therapist (MFT), or other similar licensed practice.[citation needed]
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There is stiff competition to gain acceptance into clinical psychology doctoral programs (acceptance rates of 2-5% are not uncommon). Clinical psychologists in the U.S. undergo many years of graduate training—usually five to seven years after the bachelor's degree—to gain demonstrable competence and experience. Licensure as a psychologist takes an additional one to two years post Ph.D./Psy.D. (licensure requires 3,000 hours of supervised training), depending on the state. Today in America, about half of all clinical psychology graduate students are being trained in Ph.D. programs that emphasize research and are conducted by universities—with the other half in Psy.D. programs, which have more focus on practice (similar to professional degrees for medicine and law).[6] Both types of doctoral programs (Ph.D. and Psy.D.) envision practicing clinical psychology in a research-based, scientifically valid manner, and most are accredited by the American Psychological Association.[7]
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APA accreditation[8] is very important for U.S. clinical, counseling, and school psychology programs because graduating from a non-accredited doctoral program may adversely affect employment prospects and present a hurdle for becoming licensed in some jurisdictions.[9][10][11][12]
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Doctorate (Ph.D. and Psy.D.) programs usually involve some variation on the following 5 to 7 year, 90-120 unit curriculum:
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Psychologists can be seen as practicing within two general categories of psychology: applied psychology which includes "practitioners" or "professionals", and research-orientated psychology which includes "scientists", or "scholars". The training models endorsed by the American Psychological Association (APA) require that applied psychologists be trained as both researchers and practitioners,[13] and that they possess advanced degrees.
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Psychologists typically have one of two degrees: PsyD or PhD. The PsyD program prepares the student only for clinical practice (e.g., testing, psychotherapy). Depending on the specialty (industrial/organizational, social, clinical, school, etc.), a PhD may be trained in clinical practice as well as in scientific methodology, to prepare for a career in academia or research. Both the PsyD and PhD programs prepare students to take state licensing exams.
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Within the two main categories are many further types of psychologists as reflected by the 56 professional classifications recognized by the APA,[14] including clinical, counseling, and school psychologists. Such professionals work with persons in a variety of therapeutic contexts. People often think of the discipline as involving only such clinical or counseling psychologists. While counseling and psychotherapy are common activities for psychologists, these applied fields are just two branches in the larger domain of psychology.[15] There are other classifications such as industrial, organizational and community psychologists, whose professionals mainly apply psychological research, theories, and techniques to "real-world" problems of business, industry, social benefit organizations, government,[16][17][18] and academia.
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Clinical psychologists receive training in a number of psychological therapies, including behavioral, cognitive, humanistic, existential, psychodynamic, and systemic approaches, as well as in-depth training in psychological testing, and to some extent, neuropsychological testing[19]
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Clinical psychologists can offer a range of professional services, including:[20]
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In practice, clinical psychologists might work with individuals, couples, families, or groups in a variety of settings, including private practices, hospitals, mental health organizations, schools, businesses, and non-profit agencies.
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Most clinical who engage in research and teaching do so within a college or university setting. Clinical psychologists may also choose to specialize in a particular field.
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Psychologists in the United States campaigned for legislative changes to enable specially trained psychologists to prescribe psychotropic medications. Legislation in Idaho, Iowa, Louisiana, New Mexico, and Illinois has granted those who complete an additional master's degree program in psychopharmacology permission to prescribe medications for mental and emotional disorders.[21] As of 2019[update], Louisiana is the only state where the licensing and regulation of the practice of psychology by medical psychologists (MPs) is regulated by a medical board (the Louisiana State Board of Medical Examiners) rather than a board of psychologists.[22] While other states have pursued prescriptive privileges, they have not succeeded. Similar legislation in the states of Hawaii and Oregon passed through their respective legislative bodies, but in each case the legislation was vetoed by the state's governor.[21]
|
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+
|
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In 1989, the U.S Department of Defense was directed to create the Psychopharmacology Demonstration Project (PDP). By 1997, ten psychologists were trained in psychopharmacology and granted the ability to prescribe psychiatric medications.[23]
|
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+
|
37 |
+
The practice of clinical psychology requires a license in the United States and Canada. Although each of the U.S. states is different in terms of requirements and licenses (see[24] and[25] for examples), there are three common requirements:[26]
|
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+
|
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All U.S. state, and Canada provincial, licensing boards are members of the Association of State and Provincial Psychology Boards (ASPPB) which created and maintains the Examination for Professional Practice in Psychology (EPPP). Many states require other examinations in addition to the EPPP, such as a jurisprudence (i.e. mental health law) examination or an oral examination.[26] Most states also require a certain number of continuing education credits per year in order to renew a license. Licensees can obtain this through various means, such as taking audited classes and attending approved workshops.
|
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|
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+
There are professions whose scope of practice overlaps with the practice of psychology (particularly with respect to providing psychotherapy) and for which a license is required.
|
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+
|
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+
To practice with the title of "psychologist", in almost all cases a doctorate degree is required (a PhD or PsyD in the U.S.). Normally, after the degree, the practitioner must fulfill a certain number of supervised postdoctoral hours ranging from 1,500-3,000 (usually taking one to two years), and passing the EPPP and any other state or provincial exams.[27] A professional in the U.S. must hold a graduate degree in psychology (MA, Psy.D., Ed.D., or Ph.D.), or have a state license to use the title psychologist.[28][29] Additional regulations vary from state to state.
|
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Although clinical psychologists and psychiatrists share the same fundamental aim—the alleviation of mental distress—their training, outlook, and methodologies are often different. Perhaps the most significant difference is that psychiatrists are licensed physicians, and, as such, psychiatrists are apt to use the medical model to assess mental health problems and to also employ psychotropic medications as a method of addressing mental health problems.[30]
|
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Psychologists generally do not prescribe medication, although in some jurisdictions they do have prescription privileges. In five US states (New Mexico, Louisiana, Illinois, Iowa, and Idaho), psychologists with post-doctoral clinical psychopharmacology training have been granted prescriptive authority for mental health disorders.[31][32]
|
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|
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Clinical psychologists receive extensive training in psychological test administration, scoring, interpretation, and reporting, while psychiatrists are not trained in psychological testing. In addition, psychologists (particularly those from Ph.D. programs) spend several years in graduate school being trained to conduct behavioral research; their training includes research design and advanced statistical analysis. While this training is available for physicians via dual MD/Ph.D. programs, it is not typically included in standard medical education, although psychiatrists may develop research skills during their residency or a psychiatry fellowship (post-residency). Psychologists from Psy.D. programs tend to have more training and experience in clinical practice (e.g. psychotherapy, testing) than those from Ph.D. programs.
|
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|
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Psychiatrists, as licensed physicians, have been trained more intensively in other areas, such as internal medicine and neurology, and may bring this knowledge to bear in identifying and treating medical or neurological conditions that present with primarily psychological symptoms such as depression, anxiety, or paranoia, e.g., hypothyroidism presenting with depressive symptoms, or pulmonary embolism with significant apprehension and anxiety.[33]
|
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|
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Licensed behavior analysts are licensed in five states to provide services for clients with substance abuse, developmental disabilities, and mental illness. This profession draws on the evidence base of applied behavior analysis and the philosophy of behaviorism. Behavior analysts have at least a master's degree in behavior analysis or in a mental health related discipline, as well as having taken at least five core courses in applied behavior analysis. Many behavior analysts have a doctorate. Most programs have a formalized internship program, and several programs are offered online. Most practitioners have passed the examination offered by the Behavior Analysis Certification Board.Behavior Analyst Certification Board The model licensing act for behavior analysts can be found at the Association for Behavior Analysis International's website.
|
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|
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In the United States, of 170,200 jobs for psychologists, 152,000 are employed in clinical, counseling, and school positions; 2,300 are employed in industrial-organizational positions, and 15,900 are in "all other" positions.
|
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|
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The median salary in the U.S., in 2012, for clinical, counseling, and school psychologists was US$69,280 and the median salary for organizational psychologists was US$83,580.[1][44]
|
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|
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Psychologists can work in applied or academic settings. Academic psychologists educate higher education students as well as conduct research, with graduate-level research being an important part of academic psychology. Academic positions can be tenured or non-tenured, with tenured positions being highly desirable.[45]
|
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|
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To become a psychologist, a person often completes a degree in psychology, but in other jurisdictions the course of study may be different and the activities performed may be similar to those of other professionals.[46]
|
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|
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In Australia, the psychology profession, and the use of the title "psychologist", is regulated by an Act of Parliament, the Health Practitioner Regulation (Administrative Arrangements) National Law Act 2008, following an agreement between state and territorial governments. Under this national law, registration of psychologists is administered by the Psychology Board of Australia (PsyBA).[47] Before July 2010, the professional registration of psychologists was governed by various state and territorial Psychology Registration Boards.[48] The Australian Psychology Accreditation Council (APAC) oversees education standards for the profession.
|
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The minimum requirements for general registration in psychology, including the right to use the title "psychologist", are an APAC approved four-year degree in psychology followed by either a two-year master's program or two years of practice supervised by a registered psychologist.[49][50] However, AHPRA (Australian Health Practitioner Regulation Agency) is currently in the process of phasing out the 4 + 2 internship pathway.[51] Once the 4 + 2 pathway is phased out, a master's degree or PhD will be required to become a psychologist in Australia. This is because of concerns about public safety, and to reduce the burden of training on employers.[52] There is also a '5 + 1' registration pathway, including a four-year APAC approved degree followed by one year of postgraduate study and one year of supervised practice.[53][54] Endorsement within a specific area of practice (e.g. clinical neuropsychology, clinical, community, counselling, educational and developmental, forensic, health, organisational or sport and exercise) requires additional qualifications.[55] These notations are not "specialist" titles (Western Australian psychologists could use "specialist" in their titles during a three-year transitional period from 17 October 2010 to 17 October 2013).[56][57][58]
|
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|
67 |
+
Membership with Australian Psychological Society (APS) differs from registration as a psychologist. The standard route to full membership (MAPS) of the APS usually requires four years of APAC-accredited undergraduate study, plus a master's or doctorate in psychology from an accredited institution. An alternate route is available for academics and practitioners who have gained appropriate experience and made a substantial contribution to the field of psychology.
|
68 |
+
|
69 |
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Restrictions apply to all individuals using the title "psychologist" in all states and territories of Australia. However, the terms "psychotherapist", "social worker", and "counselor" are currently self-regulated, with several organizations campaigning for government regulation.[59]
|
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|
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Since 1933, the title "psychologist" has been protected by law in Belgium. It can only be used by people who are on the National Government Commission list. The minimum requirement is the completion of five years of university training in psychology (master's degree or equivalent). The title of "psychotherapist" is not legally protected. As of 2016, Belgian law recognizes the clinical psychologist as an autonomous health profession. It reserves the practice of psychotherapy to medical doctors, clinical psychologists and clinical orthopedagogists.[60]
|
72 |
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|
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A professional in the U.S. or Canada must hold a graduate degree in psychology (MA, Psy.D., Ed.D., or Ph.D.), or have a state license to use the title psychologist.[28]
|
74 |
+
|
75 |
+
In Finland, the title "psychologist" is protected by law. The restriction for psychologists (licensed professionals) is governed by National Supervisory Authority for Welfare and Health (Finland) (Valvira).[61] It takes 330 ECTS-credits (about six years) to complete the university studies (master's degree). There are about 6,200 licensed psychologists in Finland.[62]
|
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+
|
77 |
+
In Germany, the use of the title Diplom-Psychologe (Dipl.-Psych.) is restricted by law, and a practitioner is legally required to hold the corresponding academic title, which is comparable to a M.Sc. degree and requires at least five years of training at a university. Originally, a diploma degree in psychology awarded in Germany included the subject of clinical psychology. With the Bologna-reform, this degree was replaced by a master's degree. The academic degree of Diplom-Psychologe or M.Sc. (Psychologie) does not include a psychotherapeutic qualification, which requires three to five years of additional training. The psychotherapeutic training combines in-depth theoretical knowledge with supervised patient care and self-reflection units. After having completed the training requirements, psychologists take a state-run exam, which, upon successful completion (Approbation), confers the official title of "psychological psychotherapist" (Psychologischer Psychotherapeut).[63] After many years of inter-professional political controversy, non-physician psychotherapy was given an adequate legal foundation through the creation of two new academic healthcare professions.[64]
|
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|
79 |
+
Since 1979, the title "psychologist" has been protected by law in Greece. It can only be used by people who hold a relevant license or certificate, which is issued by the Greek authorities, to practice as a psychologist. The minimum requirement is the completion of university training in psychology at a Greek university, or at a university recognized by the Greek authorities.[65] Psychologists in Greece are legally required to abide by the Code of Conduct of Psychologists (2019).[66] Psychologists in Greece are not required to register with any psychology body in the country in order to legally practice the profession.
|
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+
|
81 |
+
In New Zealand, the use of the title "psychologist" is restricted by law. Prior to 2004, only the title "registered psychologist" was restricted to people qualified and registered as such. However, with the proclamation of the Health Practitioners Competence Assurance Act, in 2003, the use of the title "psychologist" was limited to practitioners registered with the New Zealand Psychologists Board. The titles "clinical psychologist", "counseling psychologist", "educational psychologist", "intern psychologist", and "trainee psychologist" are similarly protected.[67] This is to protect the public by providing assurance that the title-holder is registered and therefore qualified and competent to practice, and can be held accountable. The legislation does not include an exemption clause for any class of practitioner (e.g., academics, or government employees).
|
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|
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In Norway, the title "psychologist" is restricted by law and can only be obtained by completing a 6 year integrated program, leading to the Candidate of Psychology degree. Psychologists are considered health personnel, and their work is regulated through the "health personnel act".[68]
|
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|
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In South Africa,[69] psychologists are qualified in either clinical, counseling, educational, organizational, or research psychology. To become qualified, one must complete a recognized master's degree in Psychology, an appropriate practicum at a recognized training institution,[70] and take an examination set by the Professional Board for Psychology.[71] Registration with the Health Professions Council of South Africa (HPCSA)[72] is required and includes a Continuing Professional Development component. The practicum usually involves a full year internship, and in some specializations, the HPCSA requires completion of an additional year of community service. The master's program consists of a seminar, coursework-based theoretical and practical training, a dissertation of limited scope, and is (in most cases) two years in duration. Prior to enrolling in the master's program, the student studies psychology for three years as an undergraduate (B.A. or B.Sc., and, for organizational psychology, also B.Com.), followed by an additional postgraduate honours degree in psychology; see List of universities in South Africa. Qualification thus requires at least five years of study and at least one internship. The undergraduate B.Psyc. is a four-year program integrating theory and practical training, and—with the required examination set by the Professional Board for Psychology—is sufficient for practice as a psychometrist or counselor.[73]
|
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|
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+
In the UK, "registered psychologist" and "practitioner psychologist" are protected titles.[74] The title of "neuropsychologist" is not protected.[74] In addition, the following specialist titles are also protected by law: "clinical psychologist", "counselling psychologist", "educational psychologist", "forensic psychologist", "health psychologist", "occupational psychologist" and "sport and exercise psychologist".[75] The Health and Care Professions Council (HCPC) is the statutory regulator for practitioner psychologists in the UK. In the UK, the use of the title "chartered psychologist" is also protected by statutory regulation, but that title simply means that the psychologist is a chartered member of the British Psychological Society, but is not necessarily registered with the HCPC. However, it is an offense for someone who is not in the appropriate section of the HCPC register to provide psychological services.[76] The requirement to register as a clinical, counseling, or educational psychologist is a professional doctorate (and in the case of the latter two the British Psychological Society's Professional Qualification, which meets the standards of a professional doctorate).[77] The title of "psychologist", by itself, is not protected.[74] The British Psychological Society is working with the HCPC to ensure that the title of "neuropsychologist" is regulated as a specialist title for practitioner psychologists.[citation needed]
|
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|
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As of December 2012[update], in the United Kingdom, there are 19,000 practitioner psychologists registered[78] across seven categories: clinical psychologist, counseling psychologist, educational psychologist, forensic psychologist, health psychologist, occupational psychologist, sport and exercise psychologist. At least 9,500 of these are clinical psychologists,[79] which is the largest group of psychologists in clinical settings such as the NHS. Around 2,000 are educational psychologists.[80]
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An avalanche (also called a snowslide) is an event that occurs when a cohesive slab of snow lying upon a weaker layer of snow fractures and slides down a steep slope. Avalanches are typically triggered in a starting zone from a mechanical failure in the snowpack (slab avalanche) when the forces of the snow exceed its strength but sometimes only with gradual widening (loose snow avalanche). After initiation, avalanches usually accelerate rapidly and grow in mass and volume as they entrain more snow. If the avalanche moves fast enough, some of the snow may mix with the air forming a powder snow avalanche, which is a type of gravity current.
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Slides of rocks or debris, behaving in a similar way to snow, are also referred to as avalanches (see rockslide[1]). The remainder of this article refers to snow avalanches, because they are the ones that occur the most in the world.
|
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The load on the snowpack may be only due to gravity, in which case failure may result either from weakening in the snowpack or increased load due to precipitation. Avalanches initiated by this process are known as spontaneous avalanches. Avalanches can also be triggered by other loading conditions such as human or biologically related activities. Seismic activity (activities in the earth's crust relating to plate movements) may also trigger the failure in the snowpack and avalanches.
|
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Although primarily composed of flowing snow and air, large avalanches have the capability to entrain ice, rocks, trees, and other surficial material. However, they are distinct from slushflows which have higher water content and more laminar flow, mudslides which have greater fluidity, rock slides which are often ice free, and serac collapses during an icefall. Avalanches are not rare or random events and are endemic to any mountain range that accumulates a standing snowpack. Avalanches are most common during winter or spring but glacier movements may cause ice and snow avalanches at any time of year. In mountainous terrain, avalanches are among the most serious objective natural hazards to life and property, with their destructive capability resulting from their potential to carry enormous masses of snow at high speeds.
|
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There is no universally accepted classification system for different forms of avalanches. Avalanches can be described by their size, their destructive potential, their initiation mechanism, their composition and their dynamics.
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Most avalanches occur spontaneously during storms under increased load due to snowfall and/or erosion. The second largest cause of natural avalanches is metamorphic changes in the snowpack such as melting due to solar radiation. Other natural causes include rain, earthquakes, rockfall and icefall. Artificial triggers of avalanches include skiers, snowmobiles, and controlled explosive work. Contrary to popular belief, avalanches are not triggered by loud sound; the pressure from sound is orders of magnitude too small to trigger an avalanche.[2]
|
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Avalanche initiation can start at a point with only a small amount of snow moving initially; this is typical of wet snow avalanches or avalanches in dry unconsolidated snow. However, if the snow has sintered into a stiff slab overlying a weak layer then fractures can propagate very rapidly, so that a large volume of snow, that may be thousands of cubic meters, can start moving almost simultaneously.
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A snowpack will fail when the load exceeds the strength. The load is straightforward; it is the weight of the snow. However, the strength of the snowpack is much more difficult to determine and is extremely heterogeneous. It varies in detail with properties of the snow grains, size, density, morphology, temperature, water content; and the properties of the bonds between the grains.
|
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[3] These properties may all metamorphose in time according to the local humidity, water vapour flux, temperature and heat flux. The top of the snowpack is also extensively influenced by incoming radiation and the local air flow. One of the aims of avalanche research is to develop and validate computer models that can describe the evolution of the seasonal snowpack over time.[4] A complicating factor is the complex interaction of terrain and weather, which causes significant spatial and temporal variability of the depths, crystal forms, and layering of the seasonal snowpack.
|
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Slab avalanches form frequently in snow that has been deposited, or redeposited by wind. They have the characteristic appearance of a block (slab) of snow cut out from its surroundings by fractures. Elements of slab avalanches include the following: a crown fracture at the top of the start zone, flank fractures on the sides of the start zones, and a fracture at the bottom called the stauchwall. The crown and flank fractures are vertical walls in the snow delineating the snow that was entrained in the avalanche from the snow that remained on the slope. Slabs can vary in thickness from a few centimetres to three metres. Slab avalanches account for around 90% of avalanche-related fatalities in backcountry users.
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The largest avalanches form turbulent suspension currents known as powder snow avalanches or mixed avalanches.[5] These consist of a powder cloud, which overlies a dense avalanche. They can form from any type of snow or initiation mechanism, but usually occur with fresh dry powder.
|
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They can exceed speeds of 300 kilometres per hour (190 mph), and masses of 10000000 tonnes; their flows can travel long distances along flat valley bottoms and even uphill for short distances.
|
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+
In contrast to powder snow avalanches, wet snow avalanches are a low velocity suspension of snow and water, with the flow confined to the track surface (McClung, first edition 1999, page 108).[3] The low speed of travel is due to the friction between the sliding surface of the track and the water saturated flow. Despite the low speed of travel (~10–40 km/h), wet snow avalanches are capable of generating powerful destructive forces, due to the large mass and density. The body of the flow of a wet snow avalanche can plough through soft snow, and can scour boulders, earth, trees, and other vegetation; leaving exposed and often scored ground in the avalanche track. Wet snow avalanches can be initiated from either loose snow releases, or slab releases, and only occur in snow packs that are water saturated and isothermally equilibrated to the melting point of water. The isothermal characteristic of wet snow avalanches has led to the secondary term of isothermal slides found in the literature (for example in Daffern, 1999, page 93).[6] At temperate latitudes wet snow avalanches are frequently associated with climatic avalanche cycles at the end of the winter season, when there is significant daytime warming.
|
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As an avalanche moves down a slope it follows a certain pathway that is dependent on the slope's degree of steepness and the volume of snow/ice involved in the mass movement. The origin of an avalanche is called the Starting Point and typically occurs on a 30–45 degree slope. The body of the pathway is called the Track of the avalanche and usually occurs on a 20–30 degree slope. When the avalanche loses its momentum and eventually stops it reaches the Runout Zone. This usually occurs when the slope has reached a steepness that is less than 20 degrees.[7] These degrees are not consistently true due to the fact that each avalanche is unique depending on the stability of the snowpack that it was derived from as well as the environmental or human influences that triggered the mass movement.
|
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People caught in avalanches can die from suffocation, trauma, or hypothermia. On average, 28 people die in avalanches every winter in the United States.[8] Globally, an average of over 150 people die each year from avalanches.
|
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+
An ice avalanche occurs when a large piece of ice, such as from a serac or calving glacier, falls onto ice (such as the Khumbu Icefall), triggering a movement of broken ice chunks. The resulting movement is more analogous to a rockfall or a landslide than a snow avalanche.[3] They are typically very difficult to predict and almost impossible to mitigate.
|
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Doug Fesler and Jill Fredston developed a conceptual model of the three primary elements of avalanches: terrain, weather, and snowpack. Terrain describes the places where avalanches occur, weather describes the meteorological conditions that create the snowpack, and snowpack describes the structural characteristics of snow that make avalanche formation possible.[3][9]
|
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Avalanche formation requires a slope shallow enough for snow to accumulate but steep enough for the snow to accelerate once set in motion by the combination of mechanical failure (of the snowpack) and gravity. The angle of the slope that can hold snow, called the angle of repose, depends on a variety of factors such as crystal form and moisture content. Some forms of drier and colder snow will only stick to shallower slopes, while wet and warm snow can bond to very steep surfaces. In particular, in coastal mountains, such as the Cordillera del Paine region of Patagonia, deep snowpacks collect on vertical and even overhanging rock faces. The slope angle that can allow moving snow to accelerate depends on a variety of factors such as the snow's shear strength (which is itself dependent upon crystal form) and the configuration of layers and inter-layer interfaces.
|
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The snowpack on slopes with sunny exposures is strongly influenced by sunshine. Diurnal cycles of thawing and refreezing can stabilize the snowpack by promoting settlement. Strong freeze-thaw cycles result in the formation of surface crusts during the night and of unstable surface snow during the day. Slopes in the lee of a ridge or of another wind obstacle accumulate more snow and are more likely to include pockets of deep snow, wind slabs, and cornices, all of which, when disturbed, may result in avalanche formation. Conversely, the snowpack on a windward slope is often much shallower than on a lee slope.
|
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Avalanches and avalanche paths share common elements: a start zone where the avalanche originates, a track along which the avalanche flows, and a runout zone where the avalanche comes to rest. The debris deposit is the accumulated mass of the avalanched snow once it has come to rest in the runout zone. For the image at left, many small avalanches form in this avalanche path every year, but most of these avalanches do not run the full vertical or horizontal length of the path. The frequency with which avalanches form in a given area is known as the return period.
|
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The start zone of an avalanche must be steep enough to allow snow to accelerate once set in motion, additionally convex slopes are less stable than concave slopes, because of the disparity between the tensile strength of snow layers and their compressive strength. The composition and structure of the ground surface beneath the snowpack influences the stability of the snowpack, either being a source of strength or weakness. Avalanches are unlikely to form in very thick forests, but boulders and sparsely distributed vegetation can create weak areas deep within the snowpack through the formation of strong temperature gradients. Full-depth avalanches (avalanches that sweep a slope virtually clean of snow cover) are more common on slopes with smooth ground, such as grass or rock slabs.
|
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|
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Generally speaking, avalanches follow drainages down-slope, frequently sharing drainage features with summertime watersheds. At and below tree line, avalanche paths through drainages are well defined by vegetation boundaries called trim lines, which occur where avalanches have removed trees and prevented regrowth of large vegetation. Engineered drainages, such as the avalanche dam on Mount Stephen in Kicking Horse Pass, have been constructed to protect people and property by redirecting the flow of avalanches. Deep debris deposits from avalanches will collect in catchments at the terminus of a run out, such as gullies and river beds.
|
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|
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+
Slopes flatter than 25 degrees or steeper than 60 degrees typically have a lower incidence of avalanches. Human-triggered avalanches have the greatest incidence when the snow's angle of repose is between 35 and 45 degrees; the critical angle, the angle at which human-triggered avalanches are most frequent, is 38 degrees. When the incidence of human triggered avalanches is normalized by the rates of recreational use, however, hazard increases uniformly with slope angle, and no significant difference in hazard for a given exposure direction can be found.[10] The rule of thumb is: A slope that is flat enough to hold snow but steep enough to ski has the potential to generate an avalanche, regardless of the angle.
|
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|
45 |
+
The snowpack is composed of ground-parallel layers that accumulate over the winter. Each layer contains ice grains that are representative of the distinct meteorological conditions during which the snow formed and was deposited. Once deposited, a snow layer continues to evolve under the influence of the meteorological conditions that prevail after deposition.
|
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|
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+
For an avalanche to occur, it is necessary that a snowpack have a weak layer (or instability) below a slab of cohesive snow. In practice the formal mechanical and structural factors related to snowpack instability are not directly observable outside of laboratories, thus the more easily observed properties of the snow layers (e.g. penetration resistance, grain size, grain type, temperature) are used as index measurements of the mechanical properties of the snow (e.g. tensile strength, friction coefficients, shear strength, and ductile strength). This results in two principal sources of uncertainty in determining snowpack stability based on snow structure: First, both the factors influencing snow stability and the specific characteristics of the snowpack vary widely within small areas and time scales, resulting in significant difficulty extrapolating point observations of snow layers across different scales of space and time. Second, the relationship between readily observable snowpack characteristics and the snowpack's critical mechanical properties has not been completely developed.
|
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|
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While the deterministic relationship between snowpack characteristics and snowpack stability is still a matter of ongoing scientific study, there is a growing empirical understanding of the snow composition and deposition characteristics that influence the likelihood of an avalanche. Observation and experience has shown that newly fallen snow requires time to bond with the snow layers beneath it, especially if the new snow falls during very cold and dry conditions. If ambient air temperatures are cold enough, shallow snow above or around boulders, plants, and other discontinuities in the slope, weakens from rapid crystal growth that occurs in the presence of a critical temperature gradient. Large, angular snow crystals are indicators of weak snow, because such crystals have fewer bonds per unit volume than small, rounded crystals that pack tightly together. Consolidated snow is less likely to slough than loose powdery layers or wet isothermal snow; however, consolidated snow is a necessary condition for the occurrence of slab avalanches, and persistent instabilities within the snowpack can hide below well-consolidated surface layers. Uncertainty associated with the empirical understanding of the factors influencing snow stability leads most professional avalanche workers to recommend conservative use of avalanche terrain relative to current snowpack instability.
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Avalanches can only occur in a standing snowpack. Typically winter seasons at high latitudes, high altitudes, or both have weather that is sufficiently unsettled and cold enough for precipitated snow to accumulate into a seasonal snowpack. Continentality, through its potentiating influence on the meteorological extremes experienced by snowpacks, is an important factor in the evolution of instabilities, and consequential occurrence of avalanches. Conversely, proximity to coastal environments moderates the meteorological extremes experienced by snowpacks, and results in a faster stabilization of the snowpack after storm cycles.[11] The evolution of the snowpack is critically sensitive to small variations within the narrow range of meteorological conditions that allow for the accumulation of snow into a snowpack. Among the critical factors controlling snowpack evolution are: heating by the sun, radiational cooling, vertical temperature gradients in standing snow, snowfall amounts, and snow types. Generally, mild winter weather will promote the settlement and stabilization of the snowpack; conversely, very cold, windy, or hot weather will weaken the snowpack.
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At temperatures close to the freezing point of water, or during times of moderate solar radiation, a gentle freeze-thaw cycle will take place. The melting and refreezing of water in the snow strengthens the snowpack during the freezing phase and weakens it during the thawing phase. A rapid rise in temperature, to a point significantly above the freezing point of water, may cause avalanche formation at any time of year.
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Persistent cold temperatures can either prevent new snow from stabilizing or destabilize the existing snowpack. Cold air temperatures on the snow surface produce a temperature gradient in the snow, because the ground temperature at the base of the snowpack is usually around 0 °C, and the ambient air temperature can be much colder. When a temperature gradient greater than 10 °C change per vertical meter of snow is sustained for more than a day, angular crystals called depth hoar or facets begin forming in the snowpack because of rapid moisture transport along the temperature gradient. These angular crystals, which bond poorly to one another and the surrounding snow, often become a persistent weakness in the snowpack. When a slab lying on top of a persistent weakness is loaded by a force greater than the strength of the slab and persistent weak layer, the persistent weak layer can fail and generate an avalanche.
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Any wind stronger than a light breeze can contribute to a rapid accumulation of snow on sheltered slopes downwind. Wind slab forms quickly and, if present, weaker snow below the slab may not have time to adjust to the new load. Even on a clear day, wind can quickly load a slope with snow by blowing snow from one place to another. Top-loading occurs when wind deposits snow from the top of a slope; cross-loading occurs when wind deposits snow parallel to the slope. When a wind blows over the top of a mountain, the leeward, or downwind, side of the mountain experiences top-loading, from the top to the bottom of that lee slope. When the wind blows across a ridge that leads up the mountain, the leeward side of the ridge is subject to cross-loading. Cross-loaded wind-slabs are usually difficult to identify visually.
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Snowstorms and rainstorms are important contributors to avalanche danger. Heavy snowfall will cause instability in the existing snowpack, both because of the additional weight and because the new snow has insufficient time to bond to underlying snow layers. Rain has a similar effect. In the short-term, rain causes instability because, like a heavy snowfall, it imposes an additional load on the snowpack; and, once rainwater seeps down through the snow, it acts as a lubricant, reducing the natural friction between snow layers that holds the snowpack together. Most avalanches happen during or soon after a storm.
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Daytime exposure to sunlight will rapidly destabilize the upper layers of the snowpack if the sunlight is strong enough to melt the snow, thereby reducing its hardness. During clear nights, the snowpack can re-freeze when ambient air temperatures fall below freezing, through the process of long-wave radiative cooling, or both. Radiative heat loss occurs when the night air is significantly cooler than the snowpack, and the heat stored in the snow is re-radiated into the atmosphere.
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When a slab avalanche forms, the slab disintegrates into increasingly smaller fragments as the snow travels downhill. If the fragments become small enough the outer layer of the avalanche, called a saltation layer, takes on the characteristics of a fluid. When sufficiently fine particles are present they can become airborne and, given a sufficient quantity of airborne snow, this portion of the avalanche can become separated from the bulk of the avalanche and travel a greater distance as a powder snow avalanche.[12] Scientific studies using radar, following the 1999 Galtür avalanche disaster, confirmed the hypothesis that a saltation layer forms between the surface and the airborne components of an avalanche, which can also separate from the bulk of the avalanche.[13]
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Driving an avalanche is the component of the avalanche's weight parallel to the slope; as the avalanche progresses any unstable snow in its path will tend to become incorporated, so increasing the overall weight. This force will increase as the steepness of the slope increases, and diminish as the slope flattens. Resisting this are a number of components that are thought to interact with each other: the friction between the avalanche and the surface beneath; friction between the air and snow within the fluid; fluid-dynamic drag at the leading edge of the avalanche; shear resistance between the avalanche and the air through which it is passing, and shear resistance between the fragments within the avalanche itself. An avalanche will continue to accelerate until the resistance exceeds the forward force.[14]
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Attempts to model avalanche behaviour date from the early 20th century, notably the work of Professor Lagotala in preparation for the 1924 Winter Olympics in Chamonix.[15] His method was developed by A. Voellmy and popularised following the publication in 1955 of his Ueber die Zerstoerungskraft von Lawinen (On the Destructive Force of Avalanches).[16]
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Voellmy used a simple empirical formula, treating an avalanche as a sliding block of snow moving with a drag force that was proportional to the square of the speed of its flow:[17]
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He and others subsequently derived other formulae that take other factors into account, with the Voellmy-Salm-Gubler and the Perla-Cheng-McClung models becoming most widely used as simple tools to model flowing (as opposed to powder snow) avalanches.[15]
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Since the 1990s many more sophisticated models have been developed. In Europe much of the recent work was carried out as part of the SATSIE (Avalanche Studies and Model Validation in Europe) research project supported by the European Commission[18] which produced the leading-edge MN2L model, now in use with the Service Restauration des Terrains en Montagne (Mountain Rescue Service) in
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France, and D2FRAM (Dynamical Two-Flow-Regime Avalanche Model), which was still undergoing validation as of 2007.[19] Other known models are the SAMOS-AT avalanche simulation software[20] and the RAMMS software.[21]
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Preventative measures are employed in areas where avalanches pose a significant threat to people, such as ski resorts, mountain towns, roads, and railways. There are several ways to prevent avalanches and lessen their power and develop preventative measures to reduce the likelihood and size of avalanches by disrupting the structure of the snowpack, while passive measures reinforce and stabilize the snowpack in situ. The simplest active measure is repeatedly traveling on a snowpack as snow accumulates; this can be by means of boot-packing, ski-cutting, or machine grooming. Explosives are used extensively to prevent avalanches, by triggering smaller avalanches that break down instabilities in the snowpack, and removing overburden that can result in larger avalanches. Explosive charges are delivered by a number of methods including hand-tossed charges, helicopter-dropped bombs, Gazex concussion lines, and ballistic projectiles launched by air cannons and artillery. Passive preventive systems such as snow fences and light walls can be used to direct the placement of snow. Snow builds up around the fence, especially the side that faces the prevailing winds. Downwind of the fence, snow buildup is lessened. This is caused by the loss of snow at the fence that would have been deposited and the pickup of the snow that is already there by the wind, which was depleted of snow at the fence. When there is a sufficient density of trees, they can greatly reduce the strength of avalanches. They hold snow in place and when there is an avalanche, the impact of the snow against the trees slows it down. Trees can either be planted or they can be conserved, such as in the building of a ski resort, to reduce the strength of avalanches.
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In turn, socio-environmental changes can influence the occurrence of damaging avalanches: some studies linking changes in land-use/land-cover patterns and the evolution of snow avalanche damage in mid latitude mountains show the importance of the role played by vegetation cover, that is at the root of the increase of damage when the protective forest is deforested (because of demographic growth, intensive grazing and industrial or legal causes), and at the root of the decrease of damage because of the transformation of a traditional land-management system based on overexploitation into a system based on land marginalization and reforestation, something that has happened mainly since the mid-20th century in mountain environments of developed countries[22]
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In many areas, regular avalanche tracks can be identified and precautions can be taken to minimise damage, such as the prevention of development in these areas. To mitigate the effect of avalanches the construction of artificial barriers can be very effective in reducing avalanche damage. There are several types: One kind of barrier (snow net) uses a net strung between poles that are anchored by guy wires in addition to their foundations. These barriers are similar to those used for rockslides. Another type of barrier is a rigid fence-like structure (snow fence) and may be constructed of steel, wood or pre-stressed concrete. They usually have gaps between the beams and are built perpendicular to the slope, with reinforcing beams on the downhill side. Rigid barriers are often considered unsightly, especially when many rows must be built. They are also expensive and vulnerable to damage from falling rocks in the warmer months. In addition to industrially manufactured barriers, landscaped barriers, called avalanche dams stop or deflect avalanches with their weight and strength. These barriers are made out of concrete, rocks or earth. They are usually placed right above the structure, road or railway that they are trying to protect, although they can also be used to channel avalanches into other barriers. Occasionally, earth mounds are placed in the avalanche's path to slow it down. Finally, along transportation corridors, large shelters, called snow sheds, can be built directly in the slide path of an avalanche to protect traffic from avalanches.
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Warning systems can detect avalanches which develop slowly, such as ice avalanches caused by icefalls from glaciers. Interferometric Radars, high-resolution Cameras, or motion sensors can monitor instable areas over a long term, lasting from days to years. Experts interpret the recorded data and are able to recognize upcoming ruptures in order to initiate appropriate measures. Such systems (e.g. the monitoring of the Weissmies glacier in Switzerland[23]) can recognize events several days in advance.
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Modern radar technology enables the monitoring of large areas and the localization of avalanches at any weather condition, by day and by night. Complex alarm systems are able to detect avalanches within a short time in order to close (e.g. roads and rails) or evacuate (e.g. construction sites) endangered areas. An example of such a system is installed on the only access road of Zermatt in Switzerland.[24] Two radars monitor the slope of a mountain above the road. The system automatically closes the road by activating several barriers and traffic lights within seconds such that no people are harmed.
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Avalanche accidents are broadly differentiated into 2 categories: accidents in recreational settings, and accidents in residential, industrial, and transportation settings. This distinction is motivated by the observed difference in the causes of avalanche accidents in the two settings. In the recreational setting most accidents are caused by the people involved in the avalanche. In a 1996 study, Jamieson et al. (pages 7–20)[25] found that 83% of all avalanches in the recreational setting were caused by those who were involved in the accident. In contrast, all of the accidents in the residential, industrial, and transportation settings were due to spontaneous natural avalanches. Because of the difference in the causes of avalanche accidents, and the activities pursued in the two settings, avalanche and disaster management professionals have developed two related preparedness, rescue, and recovery strategies for each of the settings.
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Two avalanches occurred in March 1910 in the Cascade and Selkirk Mountain ranges; On March 1 the Wellington avalanche killed 96 in Washington state, United States. Three days later 62 railroad workers were killed in the Rogers Pass avalanche in British Columbia, Canada.
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During World War I, an estimated 40,000 to 80,000 soldiers died as a result of avalanches during the mountain campaign in the Alps at the Austrian-Italian front, many of which were caused by artillery fire.[26][27] Some 10,000 men, from both sides, lost their lives in avalanches in December 1916.[28]
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In the northern hemisphere winter of 1950–1951 approximately 649 avalanches were recorded in a three-month period throughout the Alps in Austria, France, Switzerland, Italy and Germany. This series of avalanches killed around 265 people and was termed the Winter of Terror.
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A mountain climbing camp on Lenin Peak, in what is now Kyrgyzstan, was wiped out in 1990 when an earthquake triggered a large avalanche that overran the camp.[29] Forty-three climbers were killed.[30]
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In 1993, the Bayburt Üzengili avalanche killed 60 individuals in Üzengili in the province of Bayburt, Turkey.
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A large avalanche in Montroc, France, in 1999, 300,000 cubic metres of snow slid on a 30° slope, achieving a speed in the region of 100 km/h (62 mph). It killed 12 people in their chalets under 100,000 tons of snow, 5 meters (16 feet) deep. The mayor of Chamonix was convicted of second-degree murder for not evacuating the area, but received a suspended sentence.[31]
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The small Austrian village of Galtür was hit by the Galtür avalanche in 1999. The village was thought to be in a safe zone but the avalanche was exceptionally large and flowed into the village. Thirty-one people died.
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On December 1, 2000, the Glory Bowl Avalanche formed on Mt. Glory which is located within the Teton Mountain Range in Wyoming, United States. Joel Roof was snowboarding recreationally in this backcountry, bowl-shaped run and triggered the avalanche. He was carried nearly 2,000 feet to the base of the mountain and was not successfully rescued.[32]
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In Europe, the avalanche risk is widely rated on the following scale, which was adopted in April 1993 to replace the earlier non-standard national schemes. Descriptions were last updated in May 2003 to enhance uniformity.[33]
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In France, most avalanche deaths occur at risk levels 3 and 4. In Switzerland most occur at levels 2 and 3. It is thought that this may be due to national differences of interpretation when assessing the risks.[34]
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[1] Stability:
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[2] additional load:
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Gradient:
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Avalanche size:[citation needed]
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In the United States and Canada, the following avalanche danger scale is used. Descriptors vary depending on country.
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There are nine different types of avalanche problems:[35][36]
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The Canadian classification for avalanche size is based upon the consequences of the avalanche. Half sizes are commonly used.[37]
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The size of avalanches are classified using two scales; size relative to destructive force or D-scale and size relative to the avalanche path or R-scale.[38][39] Both size scales range from 1 to 5 with the D size scale half sizes can be used.[38][39]
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Slab avalanche hazard analysis can be done using the Rutschblock Test. A 2 m wide block of snow is isolated from the rest of the slope and progressively loaded. The result is a rating of slope stability on a seven step scale.[40]
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(Rutsch means slide in German).
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Climate change-caused temperature increases and changes in precipitation patterns will likely differ between the different mountain regions.[41] But overall a rising seasonal snow line and a decrease in the number of days with snow cover are predicted.[41][42] The impacts of these changes on avalanches are different at different elevations. At lower elevations a long-term reduction in the number of avalanches corresponding to a decrease in snow, and a short-term increase in the number of wet avalanches are predicted.[41][43][44][45] Higher elevations predicted to remain about the seasonal snow line, will likely see an increase in avalanche activity in regions seeing increases in precipitation during the winter season.[44][46] Predictions also show an increase in the number of rain on snow events,[42] and wet avalanche cycles occurring earlier in the spring during the remainder of this century.[47]
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Media related to Avalanche chute at Wikimedia Commons
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Claudius Ptolemy (/ˈtɒləmi/; Koinē Greek: Κλαύδιος Πτολεμαῖος, Klaúdios Ptolemaîos [kláwdios ptolɛmɛ́os]; Latin: Claudius Ptolemaeus; c. 100 – c. 170)[2] was a mathematician, astronomer, geographer and astrologer who wrote several scientific treatises, three of which were of importance to later Byzantine, Islamic and Western European science. The first is the astronomical treatise now known as the Almagest, although it was originally entitled the Mathematical Treatise (Μαθηματικὴ Σύνταξις) and then known as The Great Treatise (Ἡ Μεγάλη Σύνταξις). The second is the Geography, which is a thorough discussion of the geographic knowledge of the Greco-Roman world. The third is the astrological treatise in which he attempted to adapt horoscopic astrology to the Aristotelian natural philosophy of his day. This is sometimes known as the Apotelesmatiká (Ἀποτελεσματικά) but more commonly known as the Tetrábiblos from the Koine Greek (Τετράβιβλος) meaning "Four Books" or by the Latin Quadripartitum.
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Ptolemy lived in the city of Alexandria in the Roman province of Egypt under the rule of the Roman Empire,[3] had a Latin name (which several historians have taken to imply he was also a Roman citizen),[4] cited Greek philosophers, and used Babylonian observations and Babylonian lunar theory. The 14th century astronomer Theodore Meliteniotes gave his birthplace as the prominent Greek city Ptolemais Hermiou (Πτολεμαΐς ‘Ερμείου) in the Thebaid (Θηβᾱΐς). This attestation is quite late, however, and there is no other evidence to confirm or contradict it.[5] He died in Alexandria around 168.[6]
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Ptolemaeus (Πτολεμαῖος Ptolemaîos) is an ancient Greek personal name. It occurs once in Greek mythology and is of Homeric form.[7] It was common among the Macedonian upper class at the time of Alexander the Great and there were several of this name among Alexander's army, one of whom made himself pharaoh in 323 BCE: Ptolemy I Soter, the first pharaoh of the Ptolemaic Kingdom. All subsequent pharaohs of Egypt until Egypt became a Roman province in 30 BCE, ending the Macedonian family's rule, were also Ptolemies.[8][citation needed]
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The name Claudius is a Roman name, belonging to the gens Claudia; the peculiar multipart form of the whole name Claudius Ptolemaeus is a Roman custom, characteristic of Roman citizens. Several historians have made the deduction that this indicates that Ptolemy would have been a Roman citizen.[10] Gerald Toomer, the translator of Ptolemy's Almagest into English, suggests that citizenship was probably granted to one of Ptolemy's ancestors by either the emperor Claudius or the emperor Nero.[11]
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The 9th century Persian astronomer Abu Maʻshar presents Ptolemy as a member of Egypt's royal lineage, stating that the descendants of the Alexandrine general and Pharaoh Ptolemy I Soter, were wise "and included Ptolemy the Wise, who composed the book of the Almagest". Abu Maʻshar recorded a belief that a different member of this royal line "composed the book on astrology and attributed it to Ptolemy". We can evidence historical confusion on this point from Abu Maʿshar's subsequent remark: "It is sometimes said that the very learned man who wrote the book of astrology also wrote the book of the Almagest. The correct answer is not known."[12] Not much positive evidence is known on the subject of Ptolemy's ancestry, apart from what can be drawn from the details of his name (see above), although modern scholars have concluded that Abu Maʻshar's account is erroneous.[13] It is no longer doubted that the astronomer who wrote the Almagest also wrote the Tetrabiblos as its astrological counterpart.[14]
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Ptolemy wrote in ancient Greek and can be shown to have utilized Babylonian astronomical data.[15][16] He might have been a Roman citizen, but was ethnically either a Greek[2][17][18] or a Hellenized Egyptian.[17][19][20] He was often known in later Arabic sources as "the Upper Egyptian",[21] suggesting he may have had origins in southern Egypt.[22] Later Arabic astronomers, geographers and physicists referred to him as Baṭlumyus (Arabic: بَطْلُمْيوس).[23]
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Ptolemy's Almagest is the only surviving comprehensive ancient treatise on astronomy. Babylonian astronomers had developed arithmetical techniques for calculating astronomical phenomena; Greek astronomers such as Hipparchus had produced geometric models for calculating celestial motions. Ptolemy, however, claimed to have derived his geometrical models from selected astronomical observations by his predecessors spanning more than 800 years, though astronomers have for centuries suspected that his models' parameters were adopted independently of observations.[24] Ptolemy presented his astronomical models in convenient tables, which could be used to compute the future or past position of the planets.[25] The Almagest also contains a star catalogue, which is a version of a catalogue created by Hipparchus. Its list of forty-eight constellations is ancestral to the modern system of constellations, but unlike the modern system they did not cover the whole sky (only the sky Hipparchus could see). Across Europe, the Middle East and North Africa in the Medieval period, it was the authoritative text on astronomy, with its author becoming an almost mythical figure, called Ptolemy, King of Alexandria.[26] The Almagest was preserved, like most of extant Classical Greek science, in Arabic manuscripts (hence its familiar name). Because of its reputation, it was widely sought and was translated twice into Latin in the 12th century, once in Sicily and again in Spain.[27] Ptolemy's model, like those of his predecessors, was geocentric and was almost universally accepted until the appearance of simpler heliocentric models during the scientific revolution.
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His Planetary Hypotheses went beyond the mathematical model of the Almagest to present a physical realization of the universe as a set of nested spheres,[28] in which he used the epicycles of his planetary model to compute the dimensions of the universe. He estimated the Sun was at an average distance of 1,210 Earth radii, while the radius of the sphere of the fixed stars was 20,000 times the radius of the Earth.[29]
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Ptolemy presented a useful tool for astronomical calculations in his Handy Tables, which tabulated all the data needed to compute the positions of the Sun, Moon and planets, the rising and setting of the stars, and eclipses of the Sun and Moon. Ptolemy's Handy Tables provided the model for later astronomical tables or zījes. In the Phaseis (Risings of the Fixed Stars), Ptolemy gave a parapegma, a star calendar or almanac, based on the appearances and disappearances of stars over the course of the solar year.[30]
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Ptolemy's second main work is his Geography (also called the Geographia), a compilation of geographical coordinates of the part of the world known to the Roman Empire during his time. He relied somewhat on the work of an earlier geographer, Marinos of Tyre, and on gazetteers of the Roman and ancient Persian Empire.[citation needed] He also acknowledged ancient astronomer Hipparchus for having provided the elevation of the north celestial pole[31] for a few cities.[32]
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The first part of the Geography is a discussion of the data and of the methods he used. As with the model of the Solar System in the Almagest, Ptolemy put all this information into a grand scheme. Following Marinos, he assigned coordinates to all the places and geographic features he knew, in a grid that spanned the globe. Latitude was measured from the equator, as it is today, but Ptolemy preferred[33] to express it as climata, the length of the longest day rather than degrees of arc: the length of the midsummer day increases from 12h to 24h as one goes from the equator to the polar circle. In books 2 through 7, he used degrees and put the meridian of 0 longitude at the most western land he knew, the "Blessed Islands", often identified as the Canary Islands, as suggested by the location of the six dots labelled the "FORTUNATA" islands near the left extreme of the blue sea of Ptolemy's map here reproduced.
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Ptolemy also devised and provided instructions on how to create maps both of the whole inhabited world (oikoumenè) and of the Roman provinces. In the second part of the Geography, he provided the necessary topographic lists, and captions for the maps. His oikoumenè spanned 180 degrees of longitude from the Blessed Islands in the Atlantic Ocean to the middle of China, and about 80 degrees of latitude from Shetland to anti-Meroe (east coast of Africa); Ptolemy was well aware that he knew about only a quarter of the globe, and an erroneous extension of China southward suggests his sources did not reach all the way to the Pacific Ocean.
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The maps in surviving manuscripts of Ptolemy's Geography, however, only date from about 1300, after the text was rediscovered by Maximus Planudes. It seems likely that the topographical tables in books 2–7 are cumulative texts – texts which were altered and added to as new knowledge became available in the centuries after Ptolemy.[34] This means that information contained in different parts of the Geography is likely to be of different dates.
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Maps based on scientific principles had been made since the time of Eratosthenes, in the 3rd century BC, but Ptolemy improved map projections. It is known from a speech by Eumenius that a world map, an orbis pictus, doubtless based on the Geography, was on display in a school in Augustodunum, Gaul in the 3rd century.[35] In the 15th century, Ptolemy's Geography began to be printed with engraved maps; the earliest printed edition with engraved maps was produced in Bologna in 1477, followed quickly by a Roman edition in 1478 (Campbell, 1987). An edition printed at Ulm in 1482, including woodcut maps, was the first one printed north of the Alps. The maps look distorted when compared to modern maps, because Ptolemy's data were inaccurate. One reason is that Ptolemy estimated the size of the Earth as too small: while Eratosthenes found 700 stadia for a great circle degree on the globe, Ptolemy uses 500 stadia in the Geography. It is highly probable that these were the same stadion, since Ptolemy switched from the former scale to the latter between the Syntaxis and the Geography, and severely readjusted longitude degrees accordingly. See also Ancient Greek units of measurement and History of geodesy.
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Because Ptolemy derived many of his key latitudes from crude longest day values, his latitudes are erroneous on average by roughly a degree (2 degrees for Byzantium, 4 degrees for Carthage), though capable ancient astronomers knew their latitudes to more like a minute. (Ptolemy's own latitude was in error by 14'.) He agreed (Geography 1.4) that longitude was best determined by simultaneous observation of lunar eclipses, yet he was so out of touch with the scientists of his day that he knew of no such data more recent than 500 years before (Arbela eclipse). When switching from 700 stadia per degree to 500, he (or Marinos) expanded longitude differences between cities accordingly (a point first realized by P. Gosselin in 1790), resulting in serious over-stretching of the Earth's east-west scale in degrees, though not distance. Achieving highly precise longitude remained a problem in geography until the application of Galileo's Jovian moon method in the 18th century. It must be added that his original topographic list cannot be reconstructed: the long tables with numbers were transmitted to posterity through copies containing many scribal errors, and people have always been adding or improving the topographic data: this is a testimony to the persistent popularity of this influential work in the history of cartography.
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Ptolemy has been referred to as "a pro-astrological authority of the highest magnitude".[36] His astrological treatise, a work in four parts, is known by the Greek term Tetrabiblos, or the Latin equivalent Quadripartitum: "Four Books". Ptolemy's own title is unknown, but may have been the term found in some Greek manuscripts: Apotelesmatika, roughly meaning "Astrological Outcomes", "Effects" or "Prognostics".[37][38]
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As a source of reference, the Tetrabiblos is said to have "enjoyed almost the authority of a Bible among the astrological writers of a thousand years or more".[39] It was first translated from Arabic into Latin by Plato of Tivoli (Tiburtinus) in 1138, while he was in Spain.[40] The Tetrabiblos is an extensive and continually reprinted treatise on the ancient principles of horoscopic astrology. That it did not quite attain the unrivaled status of the Almagest was, perhaps, because it did not cover some popular areas of the subject, particularly electional astrology (interpreting astrological charts for a particular moment to determine the outcome of a course of action to be initiated at that time), and medical astrology, which were later adoptions.
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The great popularity that the Tetrabiblos did possess might be attributed to its nature as an exposition of the art of astrology, and as a compendium of astrological lore, rather than as a manual. It speaks in general terms, avoiding illustrations and details of practice. Ptolemy was concerned to defend astrology by defining its limits, compiling astronomical data that he believed was reliable and dismissing practices (such as considering the numerological significance of names) that he believed to be without sound basis.
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Much of the content of the Tetrabiblos was collected from earlier sources; Ptolemy's achievement was to order his material in a systematic way, showing how the subject could, in his view, be rationalized. It is, indeed, presented as the second part of the study of astronomy of which the Almagest was the first, concerned with the influences of the celestial bodies in the sublunary sphere. Thus explanations of a sort are provided for the astrological effects of the planets, based upon their combined effects of heating, cooling, moistening, and drying.
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Ptolemy's astrological outlook was quite practical: he thought that astrology was like medicine, that is conjectural, because of the many variable factors to be taken into account: the race, country, and upbringing of a person affects an individual's personality as much as, if not more than, the positions of the Sun, Moon, and planets at the precise moment of their birth, so Ptolemy saw astrology as something to be used in life but in no way relied on entirely.
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A collection of one hundred aphorisms about astrology called the Centiloquium, ascribed to Ptolemy, was widely reproduced and commented on by Arabic, Latin and Hebrew scholars, and often bound together in medieval manuscripts after the Tetrabiblos as a kind of summation. It is now believed to be a much later pseudepigraphical composition. The identity and date of the actual author of the work, referred to now as Pseudo-Ptolemy, remains the subject of conjecture.[dubious – discuss]
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Despite Ptolemy's prominence as a philosopher, the Dutch historian of science Eduard Jan Dijksterhuis criticizes the Tetrabiblos, stating that "it only remain puzzling that the very writer of the Almagest, who had taught how to develop astronomy from accurate observations and mathematical constructions, could put together such a system of superficial analogies and unfounded assertions."[41]
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Ptolemy also wrote an influential work, Harmonics, on music theory and the mathematics of music.[42] After criticizing the approaches of his predecessors, Ptolemy argued for basing musical intervals on mathematical ratios (in contrast to the followers of Aristoxenus and in agreement with the followers of Pythagoras), backed up by empirical observation (in contrast to the overly theoretical approach of the Pythagoreans). Ptolemy wrote about how musical notes could be translated into mathematical equations and vice versa in Harmonics. This is called Pythagorean tuning because it was first discovered by Pythagoras. However, Pythagoras believed that the mathematics of music should be based on the specific ratio of 3:2, whereas Ptolemy merely believed that it should just generally involve tetrachords and octaves. He presented his own divisions of the tetrachord and the octave, which he derived with the help of a monochord. His Harmonics never had the influence of his Almagest or Planetary Hypotheses, but a part of it (Book III) did encourage Kepler in his own musings on the harmony of the world (Kepler, Harmonice Mundi, Appendix to Book V).[43] Ptolemy's astronomical interests also appeared in a discussion of the "music of the spheres".
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His Optics is a work that survives only in a poor Arabic translation and in about twenty manuscripts of a Latin version of the Arabic, which was translated by Eugenius of Palermo (c. 1154). In it, Ptolemy writes about properties of light, including reflection, refraction, and colour. The work is a significant part of the early history of optics[44]
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and influenced the more famous 11th-century Book of Optics by Ibn al-Haytham. It contains the earliest surviving table of refraction from air to water, for which the values (with the exception of the 60° angle of incidence), although historically praised as experimentally derived, appear to have been obtained from an arithmetic progression.[45]
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The work is also important for the early history of perception. Ptolemy combined the mathematical, philosophical and physiological traditions. He held an extramission-intromission theory of vision: the rays (or flux) from the eye formed a cone, the vertex being within the eye, and the base defining the visual field. The rays were sensitive, and conveyed information back to the observer's intellect about the distance and orientation of surfaces. Size and shape were determined by the visual angle subtended at the eye combined with perceived distance and orientation. This was one of the early statements of size-distance invariance as a cause of perceptual size and shape constancy, a view supported by the Stoics.[46] Ptolemy offered explanations for many phenomena concerning illumination and colour, size, shape, movement and binocular vision. He also divided illusions into those caused by physical or optical factors and those caused by judgmental factors. He offered an obscure explanation of the sun or moon illusion (the enlarged apparent size on the horizon) based on the difficulty of looking upwards.[47][48]
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There are several characters or items named after Ptolemy, including:
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[T]he only place mentioned in any of Ptolemy's observations is Alexandria, and there is no reason to suppose that he ever lived anywhere else. The statement by Theodore Meliteniotes that he was born in Ptolemais Hermiou (in Upper Egypt) could be correct, but it is late (ca. 1360) and unsupported.
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But what we really want to know is to what extent the Alexandrian mathematicians of the period from the 1st to the 5th centuries AD were Greek. Certainly, all of them wrote in Greek and were part of the Greek intellectual community of Alexandria. Most modern studies conclude that the Greek community coexisted ... So should we assume that Ptolemy and Diophantus, Pappus and Hypatia were ethnically Greek, that their ancestors had come from Greece at some point in the past but had remained effectively isolated from the Egyptians? It is, of course, impossible to answer this question definitively. But research in papyri dating from the early centuries of the common era demonstrates that a significant amount of intermarriage took place between the Greek and Egyptian communities ... And it is known that Greek marriage contracts increasingly came to resemble Egyptian ones. In addition, even from the founding of Alexandria, small numbers of Egyptians were admitted to the privileged classes in the city to fulfill numerous civic roles. Of course, it was essential in such cases for the Egyptians to become "Hellenized", to adopt Greek habits and the Greek language. Given that the Alexandrian mathematicians mentioned here were active several hundred years after the founding of the city, it would seem at least equally possible that they were ethnically Egyptian as that they remained ethnically Greek. In any case, it is unreasonable to portray them with purely European features when no physical descriptions exist.
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Claudius Ptolemy (/ˈtɒləmi/; Koinē Greek: Κλαύδιος Πτολεμαῖος, Klaúdios Ptolemaîos [kláwdios ptolɛmɛ́os]; Latin: Claudius Ptolemaeus; c. 100 – c. 170)[2] was a mathematician, astronomer, geographer and astrologer who wrote several scientific treatises, three of which were of importance to later Byzantine, Islamic and Western European science. The first is the astronomical treatise now known as the Almagest, although it was originally entitled the Mathematical Treatise (Μαθηματικὴ Σύνταξις) and then known as The Great Treatise (Ἡ Μεγάλη Σύνταξις). The second is the Geography, which is a thorough discussion of the geographic knowledge of the Greco-Roman world. The third is the astrological treatise in which he attempted to adapt horoscopic astrology to the Aristotelian natural philosophy of his day. This is sometimes known as the Apotelesmatiká (Ἀποτελεσματικά) but more commonly known as the Tetrábiblos from the Koine Greek (Τετράβιβλος) meaning "Four Books" or by the Latin Quadripartitum.
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Ptolemy lived in the city of Alexandria in the Roman province of Egypt under the rule of the Roman Empire,[3] had a Latin name (which several historians have taken to imply he was also a Roman citizen),[4] cited Greek philosophers, and used Babylonian observations and Babylonian lunar theory. The 14th century astronomer Theodore Meliteniotes gave his birthplace as the prominent Greek city Ptolemais Hermiou (Πτολεμαΐς ‘Ερμείου) in the Thebaid (Θηβᾱΐς). This attestation is quite late, however, and there is no other evidence to confirm or contradict it.[5] He died in Alexandria around 168.[6]
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Ptolemaeus (Πτολεμαῖος Ptolemaîos) is an ancient Greek personal name. It occurs once in Greek mythology and is of Homeric form.[7] It was common among the Macedonian upper class at the time of Alexander the Great and there were several of this name among Alexander's army, one of whom made himself pharaoh in 323 BCE: Ptolemy I Soter, the first pharaoh of the Ptolemaic Kingdom. All subsequent pharaohs of Egypt until Egypt became a Roman province in 30 BCE, ending the Macedonian family's rule, were also Ptolemies.[8][citation needed]
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The name Claudius is a Roman name, belonging to the gens Claudia; the peculiar multipart form of the whole name Claudius Ptolemaeus is a Roman custom, characteristic of Roman citizens. Several historians have made the deduction that this indicates that Ptolemy would have been a Roman citizen.[10] Gerald Toomer, the translator of Ptolemy's Almagest into English, suggests that citizenship was probably granted to one of Ptolemy's ancestors by either the emperor Claudius or the emperor Nero.[11]
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The 9th century Persian astronomer Abu Maʻshar presents Ptolemy as a member of Egypt's royal lineage, stating that the descendants of the Alexandrine general and Pharaoh Ptolemy I Soter, were wise "and included Ptolemy the Wise, who composed the book of the Almagest". Abu Maʻshar recorded a belief that a different member of this royal line "composed the book on astrology and attributed it to Ptolemy". We can evidence historical confusion on this point from Abu Maʿshar's subsequent remark: "It is sometimes said that the very learned man who wrote the book of astrology also wrote the book of the Almagest. The correct answer is not known."[12] Not much positive evidence is known on the subject of Ptolemy's ancestry, apart from what can be drawn from the details of his name (see above), although modern scholars have concluded that Abu Maʻshar's account is erroneous.[13] It is no longer doubted that the astronomer who wrote the Almagest also wrote the Tetrabiblos as its astrological counterpart.[14]
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Ptolemy wrote in ancient Greek and can be shown to have utilized Babylonian astronomical data.[15][16] He might have been a Roman citizen, but was ethnically either a Greek[2][17][18] or a Hellenized Egyptian.[17][19][20] He was often known in later Arabic sources as "the Upper Egyptian",[21] suggesting he may have had origins in southern Egypt.[22] Later Arabic astronomers, geographers and physicists referred to him as Baṭlumyus (Arabic: بَطْلُمْيوس).[23]
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Ptolemy's Almagest is the only surviving comprehensive ancient treatise on astronomy. Babylonian astronomers had developed arithmetical techniques for calculating astronomical phenomena; Greek astronomers such as Hipparchus had produced geometric models for calculating celestial motions. Ptolemy, however, claimed to have derived his geometrical models from selected astronomical observations by his predecessors spanning more than 800 years, though astronomers have for centuries suspected that his models' parameters were adopted independently of observations.[24] Ptolemy presented his astronomical models in convenient tables, which could be used to compute the future or past position of the planets.[25] The Almagest also contains a star catalogue, which is a version of a catalogue created by Hipparchus. Its list of forty-eight constellations is ancestral to the modern system of constellations, but unlike the modern system they did not cover the whole sky (only the sky Hipparchus could see). Across Europe, the Middle East and North Africa in the Medieval period, it was the authoritative text on astronomy, with its author becoming an almost mythical figure, called Ptolemy, King of Alexandria.[26] The Almagest was preserved, like most of extant Classical Greek science, in Arabic manuscripts (hence its familiar name). Because of its reputation, it was widely sought and was translated twice into Latin in the 12th century, once in Sicily and again in Spain.[27] Ptolemy's model, like those of his predecessors, was geocentric and was almost universally accepted until the appearance of simpler heliocentric models during the scientific revolution.
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His Planetary Hypotheses went beyond the mathematical model of the Almagest to present a physical realization of the universe as a set of nested spheres,[28] in which he used the epicycles of his planetary model to compute the dimensions of the universe. He estimated the Sun was at an average distance of 1,210 Earth radii, while the radius of the sphere of the fixed stars was 20,000 times the radius of the Earth.[29]
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Ptolemy presented a useful tool for astronomical calculations in his Handy Tables, which tabulated all the data needed to compute the positions of the Sun, Moon and planets, the rising and setting of the stars, and eclipses of the Sun and Moon. Ptolemy's Handy Tables provided the model for later astronomical tables or zījes. In the Phaseis (Risings of the Fixed Stars), Ptolemy gave a parapegma, a star calendar or almanac, based on the appearances and disappearances of stars over the course of the solar year.[30]
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Ptolemy's second main work is his Geography (also called the Geographia), a compilation of geographical coordinates of the part of the world known to the Roman Empire during his time. He relied somewhat on the work of an earlier geographer, Marinos of Tyre, and on gazetteers of the Roman and ancient Persian Empire.[citation needed] He also acknowledged ancient astronomer Hipparchus for having provided the elevation of the north celestial pole[31] for a few cities.[32]
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The first part of the Geography is a discussion of the data and of the methods he used. As with the model of the Solar System in the Almagest, Ptolemy put all this information into a grand scheme. Following Marinos, he assigned coordinates to all the places and geographic features he knew, in a grid that spanned the globe. Latitude was measured from the equator, as it is today, but Ptolemy preferred[33] to express it as climata, the length of the longest day rather than degrees of arc: the length of the midsummer day increases from 12h to 24h as one goes from the equator to the polar circle. In books 2 through 7, he used degrees and put the meridian of 0 longitude at the most western land he knew, the "Blessed Islands", often identified as the Canary Islands, as suggested by the location of the six dots labelled the "FORTUNATA" islands near the left extreme of the blue sea of Ptolemy's map here reproduced.
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Ptolemy also devised and provided instructions on how to create maps both of the whole inhabited world (oikoumenè) and of the Roman provinces. In the second part of the Geography, he provided the necessary topographic lists, and captions for the maps. His oikoumenè spanned 180 degrees of longitude from the Blessed Islands in the Atlantic Ocean to the middle of China, and about 80 degrees of latitude from Shetland to anti-Meroe (east coast of Africa); Ptolemy was well aware that he knew about only a quarter of the globe, and an erroneous extension of China southward suggests his sources did not reach all the way to the Pacific Ocean.
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The maps in surviving manuscripts of Ptolemy's Geography, however, only date from about 1300, after the text was rediscovered by Maximus Planudes. It seems likely that the topographical tables in books 2–7 are cumulative texts – texts which were altered and added to as new knowledge became available in the centuries after Ptolemy.[34] This means that information contained in different parts of the Geography is likely to be of different dates.
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Maps based on scientific principles had been made since the time of Eratosthenes, in the 3rd century BC, but Ptolemy improved map projections. It is known from a speech by Eumenius that a world map, an orbis pictus, doubtless based on the Geography, was on display in a school in Augustodunum, Gaul in the 3rd century.[35] In the 15th century, Ptolemy's Geography began to be printed with engraved maps; the earliest printed edition with engraved maps was produced in Bologna in 1477, followed quickly by a Roman edition in 1478 (Campbell, 1987). An edition printed at Ulm in 1482, including woodcut maps, was the first one printed north of the Alps. The maps look distorted when compared to modern maps, because Ptolemy's data were inaccurate. One reason is that Ptolemy estimated the size of the Earth as too small: while Eratosthenes found 700 stadia for a great circle degree on the globe, Ptolemy uses 500 stadia in the Geography. It is highly probable that these were the same stadion, since Ptolemy switched from the former scale to the latter between the Syntaxis and the Geography, and severely readjusted longitude degrees accordingly. See also Ancient Greek units of measurement and History of geodesy.
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Because Ptolemy derived many of his key latitudes from crude longest day values, his latitudes are erroneous on average by roughly a degree (2 degrees for Byzantium, 4 degrees for Carthage), though capable ancient astronomers knew their latitudes to more like a minute. (Ptolemy's own latitude was in error by 14'.) He agreed (Geography 1.4) that longitude was best determined by simultaneous observation of lunar eclipses, yet he was so out of touch with the scientists of his day that he knew of no such data more recent than 500 years before (Arbela eclipse). When switching from 700 stadia per degree to 500, he (or Marinos) expanded longitude differences between cities accordingly (a point first realized by P. Gosselin in 1790), resulting in serious over-stretching of the Earth's east-west scale in degrees, though not distance. Achieving highly precise longitude remained a problem in geography until the application of Galileo's Jovian moon method in the 18th century. It must be added that his original topographic list cannot be reconstructed: the long tables with numbers were transmitted to posterity through copies containing many scribal errors, and people have always been adding or improving the topographic data: this is a testimony to the persistent popularity of this influential work in the history of cartography.
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Ptolemy has been referred to as "a pro-astrological authority of the highest magnitude".[36] His astrological treatise, a work in four parts, is known by the Greek term Tetrabiblos, or the Latin equivalent Quadripartitum: "Four Books". Ptolemy's own title is unknown, but may have been the term found in some Greek manuscripts: Apotelesmatika, roughly meaning "Astrological Outcomes", "Effects" or "Prognostics".[37][38]
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As a source of reference, the Tetrabiblos is said to have "enjoyed almost the authority of a Bible among the astrological writers of a thousand years or more".[39] It was first translated from Arabic into Latin by Plato of Tivoli (Tiburtinus) in 1138, while he was in Spain.[40] The Tetrabiblos is an extensive and continually reprinted treatise on the ancient principles of horoscopic astrology. That it did not quite attain the unrivaled status of the Almagest was, perhaps, because it did not cover some popular areas of the subject, particularly electional astrology (interpreting astrological charts for a particular moment to determine the outcome of a course of action to be initiated at that time), and medical astrology, which were later adoptions.
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The great popularity that the Tetrabiblos did possess might be attributed to its nature as an exposition of the art of astrology, and as a compendium of astrological lore, rather than as a manual. It speaks in general terms, avoiding illustrations and details of practice. Ptolemy was concerned to defend astrology by defining its limits, compiling astronomical data that he believed was reliable and dismissing practices (such as considering the numerological significance of names) that he believed to be without sound basis.
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Much of the content of the Tetrabiblos was collected from earlier sources; Ptolemy's achievement was to order his material in a systematic way, showing how the subject could, in his view, be rationalized. It is, indeed, presented as the second part of the study of astronomy of which the Almagest was the first, concerned with the influences of the celestial bodies in the sublunary sphere. Thus explanations of a sort are provided for the astrological effects of the planets, based upon their combined effects of heating, cooling, moistening, and drying.
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Ptolemy's astrological outlook was quite practical: he thought that astrology was like medicine, that is conjectural, because of the many variable factors to be taken into account: the race, country, and upbringing of a person affects an individual's personality as much as, if not more than, the positions of the Sun, Moon, and planets at the precise moment of their birth, so Ptolemy saw astrology as something to be used in life but in no way relied on entirely.
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A collection of one hundred aphorisms about astrology called the Centiloquium, ascribed to Ptolemy, was widely reproduced and commented on by Arabic, Latin and Hebrew scholars, and often bound together in medieval manuscripts after the Tetrabiblos as a kind of summation. It is now believed to be a much later pseudepigraphical composition. The identity and date of the actual author of the work, referred to now as Pseudo-Ptolemy, remains the subject of conjecture.[dubious – discuss]
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Despite Ptolemy's prominence as a philosopher, the Dutch historian of science Eduard Jan Dijksterhuis criticizes the Tetrabiblos, stating that "it only remain puzzling that the very writer of the Almagest, who had taught how to develop astronomy from accurate observations and mathematical constructions, could put together such a system of superficial analogies and unfounded assertions."[41]
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Ptolemy also wrote an influential work, Harmonics, on music theory and the mathematics of music.[42] After criticizing the approaches of his predecessors, Ptolemy argued for basing musical intervals on mathematical ratios (in contrast to the followers of Aristoxenus and in agreement with the followers of Pythagoras), backed up by empirical observation (in contrast to the overly theoretical approach of the Pythagoreans). Ptolemy wrote about how musical notes could be translated into mathematical equations and vice versa in Harmonics. This is called Pythagorean tuning because it was first discovered by Pythagoras. However, Pythagoras believed that the mathematics of music should be based on the specific ratio of 3:2, whereas Ptolemy merely believed that it should just generally involve tetrachords and octaves. He presented his own divisions of the tetrachord and the octave, which he derived with the help of a monochord. His Harmonics never had the influence of his Almagest or Planetary Hypotheses, but a part of it (Book III) did encourage Kepler in his own musings on the harmony of the world (Kepler, Harmonice Mundi, Appendix to Book V).[43] Ptolemy's astronomical interests also appeared in a discussion of the "music of the spheres".
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+
His Optics is a work that survives only in a poor Arabic translation and in about twenty manuscripts of a Latin version of the Arabic, which was translated by Eugenius of Palermo (c. 1154). In it, Ptolemy writes about properties of light, including reflection, refraction, and colour. The work is a significant part of the early history of optics[44]
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+
and influenced the more famous 11th-century Book of Optics by Ibn al-Haytham. It contains the earliest surviving table of refraction from air to water, for which the values (with the exception of the 60° angle of incidence), although historically praised as experimentally derived, appear to have been obtained from an arithmetic progression.[45]
|
51 |
+
|
52 |
+
The work is also important for the early history of perception. Ptolemy combined the mathematical, philosophical and physiological traditions. He held an extramission-intromission theory of vision: the rays (or flux) from the eye formed a cone, the vertex being within the eye, and the base defining the visual field. The rays were sensitive, and conveyed information back to the observer's intellect about the distance and orientation of surfaces. Size and shape were determined by the visual angle subtended at the eye combined with perceived distance and orientation. This was one of the early statements of size-distance invariance as a cause of perceptual size and shape constancy, a view supported by the Stoics.[46] Ptolemy offered explanations for many phenomena concerning illumination and colour, size, shape, movement and binocular vision. He also divided illusions into those caused by physical or optical factors and those caused by judgmental factors. He offered an obscure explanation of the sun or moon illusion (the enlarged apparent size on the horizon) based on the difficulty of looking upwards.[47][48]
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There are several characters or items named after Ptolemy, including:
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[T]he only place mentioned in any of Ptolemy's observations is Alexandria, and there is no reason to suppose that he ever lived anywhere else. The statement by Theodore Meliteniotes that he was born in Ptolemais Hermiou (in Upper Egypt) could be correct, but it is late (ca. 1360) and unsupported.
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But what we really want to know is to what extent the Alexandrian mathematicians of the period from the 1st to the 5th centuries AD were Greek. Certainly, all of them wrote in Greek and were part of the Greek intellectual community of Alexandria. Most modern studies conclude that the Greek community coexisted ... So should we assume that Ptolemy and Diophantus, Pappus and Hypatia were ethnically Greek, that their ancestors had come from Greece at some point in the past but had remained effectively isolated from the Egyptians? It is, of course, impossible to answer this question definitively. But research in papyri dating from the early centuries of the common era demonstrates that a significant amount of intermarriage took place between the Greek and Egyptian communities ... And it is known that Greek marriage contracts increasingly came to resemble Egyptian ones. In addition, even from the founding of Alexandria, small numbers of Egyptians were admitted to the privileged classes in the city to fulfill numerous civic roles. Of course, it was essential in such cases for the Egyptians to become "Hellenized", to adopt Greek habits and the Greek language. Given that the Alexandrian mathematicians mentioned here were active several hundred years after the founding of the city, it would seem at least equally possible that they were ethnically Egyptian as that they remained ethnically Greek. In any case, it is unreasonable to portray them with purely European features when no physical descriptions exist.
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Claudius Ptolemy (/ˈtɒləmi/; Koinē Greek: Κλαύδιος Πτολεμαῖος, Klaúdios Ptolemaîos [kláwdios ptolɛmɛ́os]; Latin: Claudius Ptolemaeus; c. 100 – c. 170)[2] was a mathematician, astronomer, geographer and astrologer who wrote several scientific treatises, three of which were of importance to later Byzantine, Islamic and Western European science. The first is the astronomical treatise now known as the Almagest, although it was originally entitled the Mathematical Treatise (Μαθηματικὴ Σύνταξις) and then known as The Great Treatise (Ἡ Μεγάλη Σύνταξις). The second is the Geography, which is a thorough discussion of the geographic knowledge of the Greco-Roman world. The third is the astrological treatise in which he attempted to adapt horoscopic astrology to the Aristotelian natural philosophy of his day. This is sometimes known as the Apotelesmatiká (Ἀποτελεσματικά) but more commonly known as the Tetrábiblos from the Koine Greek (Τετράβιβλος) meaning "Four Books" or by the Latin Quadripartitum.
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Ptolemy lived in the city of Alexandria in the Roman province of Egypt under the rule of the Roman Empire,[3] had a Latin name (which several historians have taken to imply he was also a Roman citizen),[4] cited Greek philosophers, and used Babylonian observations and Babylonian lunar theory. The 14th century astronomer Theodore Meliteniotes gave his birthplace as the prominent Greek city Ptolemais Hermiou (Πτολεμαΐς ‘Ερμείου) in the Thebaid (Θηβᾱΐς). This attestation is quite late, however, and there is no other evidence to confirm or contradict it.[5] He died in Alexandria around 168.[6]
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Ptolemaeus (Πτολεμαῖος Ptolemaîos) is an ancient Greek personal name. It occurs once in Greek mythology and is of Homeric form.[7] It was common among the Macedonian upper class at the time of Alexander the Great and there were several of this name among Alexander's army, one of whom made himself pharaoh in 323 BCE: Ptolemy I Soter, the first pharaoh of the Ptolemaic Kingdom. All subsequent pharaohs of Egypt until Egypt became a Roman province in 30 BCE, ending the Macedonian family's rule, were also Ptolemies.[8][citation needed]
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The name Claudius is a Roman name, belonging to the gens Claudia; the peculiar multipart form of the whole name Claudius Ptolemaeus is a Roman custom, characteristic of Roman citizens. Several historians have made the deduction that this indicates that Ptolemy would have been a Roman citizen.[10] Gerald Toomer, the translator of Ptolemy's Almagest into English, suggests that citizenship was probably granted to one of Ptolemy's ancestors by either the emperor Claudius or the emperor Nero.[11]
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The 9th century Persian astronomer Abu Maʻshar presents Ptolemy as a member of Egypt's royal lineage, stating that the descendants of the Alexandrine general and Pharaoh Ptolemy I Soter, were wise "and included Ptolemy the Wise, who composed the book of the Almagest". Abu Maʻshar recorded a belief that a different member of this royal line "composed the book on astrology and attributed it to Ptolemy". We can evidence historical confusion on this point from Abu Maʿshar's subsequent remark: "It is sometimes said that the very learned man who wrote the book of astrology also wrote the book of the Almagest. The correct answer is not known."[12] Not much positive evidence is known on the subject of Ptolemy's ancestry, apart from what can be drawn from the details of his name (see above), although modern scholars have concluded that Abu Maʻshar's account is erroneous.[13] It is no longer doubted that the astronomer who wrote the Almagest also wrote the Tetrabiblos as its astrological counterpart.[14]
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Ptolemy wrote in ancient Greek and can be shown to have utilized Babylonian astronomical data.[15][16] He might have been a Roman citizen, but was ethnically either a Greek[2][17][18] or a Hellenized Egyptian.[17][19][20] He was often known in later Arabic sources as "the Upper Egyptian",[21] suggesting he may have had origins in southern Egypt.[22] Later Arabic astronomers, geographers and physicists referred to him as Baṭlumyus (Arabic: بَطْلُمْيوس).[23]
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Ptolemy's Almagest is the only surviving comprehensive ancient treatise on astronomy. Babylonian astronomers had developed arithmetical techniques for calculating astronomical phenomena; Greek astronomers such as Hipparchus had produced geometric models for calculating celestial motions. Ptolemy, however, claimed to have derived his geometrical models from selected astronomical observations by his predecessors spanning more than 800 years, though astronomers have for centuries suspected that his models' parameters were adopted independently of observations.[24] Ptolemy presented his astronomical models in convenient tables, which could be used to compute the future or past position of the planets.[25] The Almagest also contains a star catalogue, which is a version of a catalogue created by Hipparchus. Its list of forty-eight constellations is ancestral to the modern system of constellations, but unlike the modern system they did not cover the whole sky (only the sky Hipparchus could see). Across Europe, the Middle East and North Africa in the Medieval period, it was the authoritative text on astronomy, with its author becoming an almost mythical figure, called Ptolemy, King of Alexandria.[26] The Almagest was preserved, like most of extant Classical Greek science, in Arabic manuscripts (hence its familiar name). Because of its reputation, it was widely sought and was translated twice into Latin in the 12th century, once in Sicily and again in Spain.[27] Ptolemy's model, like those of his predecessors, was geocentric and was almost universally accepted until the appearance of simpler heliocentric models during the scientific revolution.
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His Planetary Hypotheses went beyond the mathematical model of the Almagest to present a physical realization of the universe as a set of nested spheres,[28] in which he used the epicycles of his planetary model to compute the dimensions of the universe. He estimated the Sun was at an average distance of 1,210 Earth radii, while the radius of the sphere of the fixed stars was 20,000 times the radius of the Earth.[29]
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Ptolemy presented a useful tool for astronomical calculations in his Handy Tables, which tabulated all the data needed to compute the positions of the Sun, Moon and planets, the rising and setting of the stars, and eclipses of the Sun and Moon. Ptolemy's Handy Tables provided the model for later astronomical tables or zījes. In the Phaseis (Risings of the Fixed Stars), Ptolemy gave a parapegma, a star calendar or almanac, based on the appearances and disappearances of stars over the course of the solar year.[30]
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Ptolemy's second main work is his Geography (also called the Geographia), a compilation of geographical coordinates of the part of the world known to the Roman Empire during his time. He relied somewhat on the work of an earlier geographer, Marinos of Tyre, and on gazetteers of the Roman and ancient Persian Empire.[citation needed] He also acknowledged ancient astronomer Hipparchus for having provided the elevation of the north celestial pole[31] for a few cities.[32]
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The first part of the Geography is a discussion of the data and of the methods he used. As with the model of the Solar System in the Almagest, Ptolemy put all this information into a grand scheme. Following Marinos, he assigned coordinates to all the places and geographic features he knew, in a grid that spanned the globe. Latitude was measured from the equator, as it is today, but Ptolemy preferred[33] to express it as climata, the length of the longest day rather than degrees of arc: the length of the midsummer day increases from 12h to 24h as one goes from the equator to the polar circle. In books 2 through 7, he used degrees and put the meridian of 0 longitude at the most western land he knew, the "Blessed Islands", often identified as the Canary Islands, as suggested by the location of the six dots labelled the "FORTUNATA" islands near the left extreme of the blue sea of Ptolemy's map here reproduced.
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Ptolemy also devised and provided instructions on how to create maps both of the whole inhabited world (oikoumenè) and of the Roman provinces. In the second part of the Geography, he provided the necessary topographic lists, and captions for the maps. His oikoumenè spanned 180 degrees of longitude from the Blessed Islands in the Atlantic Ocean to the middle of China, and about 80 degrees of latitude from Shetland to anti-Meroe (east coast of Africa); Ptolemy was well aware that he knew about only a quarter of the globe, and an erroneous extension of China southward suggests his sources did not reach all the way to the Pacific Ocean.
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The maps in surviving manuscripts of Ptolemy's Geography, however, only date from about 1300, after the text was rediscovered by Maximus Planudes. It seems likely that the topographical tables in books 2–7 are cumulative texts – texts which were altered and added to as new knowledge became available in the centuries after Ptolemy.[34] This means that information contained in different parts of the Geography is likely to be of different dates.
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Maps based on scientific principles had been made since the time of Eratosthenes, in the 3rd century BC, but Ptolemy improved map projections. It is known from a speech by Eumenius that a world map, an orbis pictus, doubtless based on the Geography, was on display in a school in Augustodunum, Gaul in the 3rd century.[35] In the 15th century, Ptolemy's Geography began to be printed with engraved maps; the earliest printed edition with engraved maps was produced in Bologna in 1477, followed quickly by a Roman edition in 1478 (Campbell, 1987). An edition printed at Ulm in 1482, including woodcut maps, was the first one printed north of the Alps. The maps look distorted when compared to modern maps, because Ptolemy's data were inaccurate. One reason is that Ptolemy estimated the size of the Earth as too small: while Eratosthenes found 700 stadia for a great circle degree on the globe, Ptolemy uses 500 stadia in the Geography. It is highly probable that these were the same stadion, since Ptolemy switched from the former scale to the latter between the Syntaxis and the Geography, and severely readjusted longitude degrees accordingly. See also Ancient Greek units of measurement and History of geodesy.
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Because Ptolemy derived many of his key latitudes from crude longest day values, his latitudes are erroneous on average by roughly a degree (2 degrees for Byzantium, 4 degrees for Carthage), though capable ancient astronomers knew their latitudes to more like a minute. (Ptolemy's own latitude was in error by 14'.) He agreed (Geography 1.4) that longitude was best determined by simultaneous observation of lunar eclipses, yet he was so out of touch with the scientists of his day that he knew of no such data more recent than 500 years before (Arbela eclipse). When switching from 700 stadia per degree to 500, he (or Marinos) expanded longitude differences between cities accordingly (a point first realized by P. Gosselin in 1790), resulting in serious over-stretching of the Earth's east-west scale in degrees, though not distance. Achieving highly precise longitude remained a problem in geography until the application of Galileo's Jovian moon method in the 18th century. It must be added that his original topographic list cannot be reconstructed: the long tables with numbers were transmitted to posterity through copies containing many scribal errors, and people have always been adding or improving the topographic data: this is a testimony to the persistent popularity of this influential work in the history of cartography.
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Ptolemy has been referred to as "a pro-astrological authority of the highest magnitude".[36] His astrological treatise, a work in four parts, is known by the Greek term Tetrabiblos, or the Latin equivalent Quadripartitum: "Four Books". Ptolemy's own title is unknown, but may have been the term found in some Greek manuscripts: Apotelesmatika, roughly meaning "Astrological Outcomes", "Effects" or "Prognostics".[37][38]
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As a source of reference, the Tetrabiblos is said to have "enjoyed almost the authority of a Bible among the astrological writers of a thousand years or more".[39] It was first translated from Arabic into Latin by Plato of Tivoli (Tiburtinus) in 1138, while he was in Spain.[40] The Tetrabiblos is an extensive and continually reprinted treatise on the ancient principles of horoscopic astrology. That it did not quite attain the unrivaled status of the Almagest was, perhaps, because it did not cover some popular areas of the subject, particularly electional astrology (interpreting astrological charts for a particular moment to determine the outcome of a course of action to be initiated at that time), and medical astrology, which were later adoptions.
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The great popularity that the Tetrabiblos did possess might be attributed to its nature as an exposition of the art of astrology, and as a compendium of astrological lore, rather than as a manual. It speaks in general terms, avoiding illustrations and details of practice. Ptolemy was concerned to defend astrology by defining its limits, compiling astronomical data that he believed was reliable and dismissing practices (such as considering the numerological significance of names) that he believed to be without sound basis.
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Much of the content of the Tetrabiblos was collected from earlier sources; Ptolemy's achievement was to order his material in a systematic way, showing how the subject could, in his view, be rationalized. It is, indeed, presented as the second part of the study of astronomy of which the Almagest was the first, concerned with the influences of the celestial bodies in the sublunary sphere. Thus explanations of a sort are provided for the astrological effects of the planets, based upon their combined effects of heating, cooling, moistening, and drying.
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Ptolemy's astrological outlook was quite practical: he thought that astrology was like medicine, that is conjectural, because of the many variable factors to be taken into account: the race, country, and upbringing of a person affects an individual's personality as much as, if not more than, the positions of the Sun, Moon, and planets at the precise moment of their birth, so Ptolemy saw astrology as something to be used in life but in no way relied on entirely.
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A collection of one hundred aphorisms about astrology called the Centiloquium, ascribed to Ptolemy, was widely reproduced and commented on by Arabic, Latin and Hebrew scholars, and often bound together in medieval manuscripts after the Tetrabiblos as a kind of summation. It is now believed to be a much later pseudepigraphical composition. The identity and date of the actual author of the work, referred to now as Pseudo-Ptolemy, remains the subject of conjecture.[dubious – discuss]
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Despite Ptolemy's prominence as a philosopher, the Dutch historian of science Eduard Jan Dijksterhuis criticizes the Tetrabiblos, stating that "it only remain puzzling that the very writer of the Almagest, who had taught how to develop astronomy from accurate observations and mathematical constructions, could put together such a system of superficial analogies and unfounded assertions."[41]
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Ptolemy also wrote an influential work, Harmonics, on music theory and the mathematics of music.[42] After criticizing the approaches of his predecessors, Ptolemy argued for basing musical intervals on mathematical ratios (in contrast to the followers of Aristoxenus and in agreement with the followers of Pythagoras), backed up by empirical observation (in contrast to the overly theoretical approach of the Pythagoreans). Ptolemy wrote about how musical notes could be translated into mathematical equations and vice versa in Harmonics. This is called Pythagorean tuning because it was first discovered by Pythagoras. However, Pythagoras believed that the mathematics of music should be based on the specific ratio of 3:2, whereas Ptolemy merely believed that it should just generally involve tetrachords and octaves. He presented his own divisions of the tetrachord and the octave, which he derived with the help of a monochord. His Harmonics never had the influence of his Almagest or Planetary Hypotheses, but a part of it (Book III) did encourage Kepler in his own musings on the harmony of the world (Kepler, Harmonice Mundi, Appendix to Book V).[43] Ptolemy's astronomical interests also appeared in a discussion of the "music of the spheres".
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His Optics is a work that survives only in a poor Arabic translation and in about twenty manuscripts of a Latin version of the Arabic, which was translated by Eugenius of Palermo (c. 1154). In it, Ptolemy writes about properties of light, including reflection, refraction, and colour. The work is a significant part of the early history of optics[44]
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and influenced the more famous 11th-century Book of Optics by Ibn al-Haytham. It contains the earliest surviving table of refraction from air to water, for which the values (with the exception of the 60° angle of incidence), although historically praised as experimentally derived, appear to have been obtained from an arithmetic progression.[45]
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The work is also important for the early history of perception. Ptolemy combined the mathematical, philosophical and physiological traditions. He held an extramission-intromission theory of vision: the rays (or flux) from the eye formed a cone, the vertex being within the eye, and the base defining the visual field. The rays were sensitive, and conveyed information back to the observer's intellect about the distance and orientation of surfaces. Size and shape were determined by the visual angle subtended at the eye combined with perceived distance and orientation. This was one of the early statements of size-distance invariance as a cause of perceptual size and shape constancy, a view supported by the Stoics.[46] Ptolemy offered explanations for many phenomena concerning illumination and colour, size, shape, movement and binocular vision. He also divided illusions into those caused by physical or optical factors and those caused by judgmental factors. He offered an obscure explanation of the sun or moon illusion (the enlarged apparent size on the horizon) based on the difficulty of looking upwards.[47][48]
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There are several characters or items named after Ptolemy, including:
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[T]he only place mentioned in any of Ptolemy's observations is Alexandria, and there is no reason to suppose that he ever lived anywhere else. The statement by Theodore Meliteniotes that he was born in Ptolemais Hermiou (in Upper Egypt) could be correct, but it is late (ca. 1360) and unsupported.
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But what we really want to know is to what extent the Alexandrian mathematicians of the period from the 1st to the 5th centuries AD were Greek. Certainly, all of them wrote in Greek and were part of the Greek intellectual community of Alexandria. Most modern studies conclude that the Greek community coexisted ... So should we assume that Ptolemy and Diophantus, Pappus and Hypatia were ethnically Greek, that their ancestors had come from Greece at some point in the past but had remained effectively isolated from the Egyptians? It is, of course, impossible to answer this question definitively. But research in papyri dating from the early centuries of the common era demonstrates that a significant amount of intermarriage took place between the Greek and Egyptian communities ... And it is known that Greek marriage contracts increasingly came to resemble Egyptian ones. In addition, even from the founding of Alexandria, small numbers of Egyptians were admitted to the privileged classes in the city to fulfill numerous civic roles. Of course, it was essential in such cases for the Egyptians to become "Hellenized", to adopt Greek habits and the Greek language. Given that the Alexandrian mathematicians mentioned here were active several hundred years after the founding of the city, it would seem at least equally possible that they were ethnically Egyptian as that they remained ethnically Greek. In any case, it is unreasonable to portray them with purely European features when no physical descriptions exist.
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Puberty is the process of physical changes through which a child's body matures into an adult body capable of sexual reproduction. It is initiated by hormonal signals from the brain to the gonads: the ovaries in a girl, the testes in a boy. In response to the signals, the gonads produce hormones that stimulate libido and the growth, function, and transformation of the brain, bones, muscle, blood, skin, hair, breasts, and sex organs. Physical growth—height and weight—accelerates in the first half of puberty and is completed when an adult body has been developed. Before puberty, the external sex organs, known as primary sexual characteristics, are sex characteristics that distinguish boys and girls. Puberty leads to sexual dimorphism through the development of the secondary sex characteristics, which further distinguish the sexes.
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On average, girls begin puberty at ages 10–11 and complete puberty at ages 15–17; boys generally begin puberty at ages 11–12 and complete puberty at ages 16–17.[1][2][3] The major landmark of puberty for females is menarche, the onset of menstruation, which occurs on average between ages 12 and 13.[2] For males, first ejaculation occurs on average at age 13.[4] In the 21st century, the average age at which children, especially girls, reach puberty is lower compared to the 19th century, when it was 15 for girls and 16 for boys.[5] This can be due to any number of factors, including improved nutrition resulting in rapid body growth, increased weight and fat deposition,[6] or exposure to endocrine disruptors such as xenoestrogens, which can at times be due to food consumption or other environmental factors.[7][8] Puberty which starts earlier than usual is known as precocious puberty, and puberty which starts later than usual is known as delayed puberty.
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Notable among the morphologic changes in size, shape, composition, and functioning of the pubertal body, is the development of secondary sex characteristics, the "filling in" of the child's body; from girl to woman, from boy to man. Derived from the Latin puberatum (age of maturity), the word puberty describes the physical changes to sexual maturation, not the psychosocial and cultural maturation denoted by the term adolescent development in Western culture, wherein adolescence is the period of mental transition from childhood to adulthood, which overlaps much of the body's period of puberty.[9]
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Two of the most significant differences between puberty in girls and puberty in boys are the age at which it begins, and the major sex steroids involved, the androgens and the estrogens.
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Although there is a wide range of normal ages, girls typically begin puberty around ages 10–11 and end puberty around 15–17; boys begin around ages 11–12 and end around 16–17.[1][2][3] Girls attain reproductive maturity about four years after the first physical changes of puberty appear.[10] In contrast, boys accelerate more slowly but continue to grow for about six years after the first visible pubertal changes.[11] Any increase in height beyond the post-pubertal age is uncommon.
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For boys, the androgen testosterone is the principal sex hormone; while testosterone is produced, all boys' changes are characterized as virilization. A substantial product of testosterone metabolism in males is estradiol. The conversion of testosterone to estradiol depends on the amount of body fat and estradiol levels in boys are typically much lower than in girls. The male "growth spurt" also begins later, accelerates more slowly, and lasts longer before the epiphyses fuse. Although boys are on average 2 centimetres (0.8 in) shorter than girls before puberty begins, adult men are on average about 13 centimetres (5.1 in) taller than women. Most of this sex difference in adult heights is attributable to a later onset of the growth spurt and a slower progression to completion, a direct result of the later rise and lower adult male levels of estradiol.[12]
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The hormone that dominates female development is an estrogen called estradiol. While estradiol promotes growth of the breasts and uterus, it is also the principal hormone driving the pubertal growth spurt and epiphyseal maturation and closure.[13] Estradiol levels rise earlier and reach higher levels in women than in men.
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The hormonal maturation of females is considerably more complicated than in boys. The main steroid hormones, testosterone, estradiol, and progesterone as well as prolactin play important physiological functions in puberty.
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Gonadal steroidgenesis in girls starts with production of testosterone which is typically quickly converted to estradiol inside the ovaries. However the rate of conversion from testosterone to estradiol (driven by FSH/LH balance) during early puberty is highly individual, resulting in very diverse development patterns of secondary sexual characteristics. Production of progesterone in the ovaries begins with the development of ovulatory cycles in girls (during the lutheal phase of the cycle), before puberty low levels of progesterone are produced in the adrenal glands of both boys and girls.
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Puberty is preceded by adrenarche, marking an increase of adrenal androgen production between ages 6–10. Adrenarche is sometimes accompanied by the early appearance of axillary and pubic hair. The first androgenic hair resulting from adrenarche can be also transient and disappear before the onset of true puberty.
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The onset of puberty is associated with high GnRH pulsing, which precedes the rise in sex hormones, LH and FSH.[14] Exogenous GnRH pulses cause the onset of puberty.[15] Brain tumors which increase GnRH output may also lead to premature puberty.[16]
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The cause of the GnRH rise is unknown. Leptin might be the cause of the GnRH rise. Leptin has receptors in the hypothalamus which synthesizes GnRH.[17] Individuals who are deficient in leptin fail to initiate puberty.[18] The levels of leptin increase with the onset of puberty, and then decline to adult levels when puberty is completed. The rise in GnRH might also be caused by genetics. A study[19] discovered that a mutation in genes encoding both neurokinin B as well as the neurokinin B receptor can alter the timing of puberty. The researchers hypothesized that neurokinin B might play a role in regulating the secretion of kisspeptin, a compound responsible for triggering direct release of GnRH as well as indirect release of LH and FSH.
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Several studies about puberty have examined the effects of an early or a late onset of puberty in males and females. In general, girls who enter puberty late experience positive outcomes in adolescence and adulthood, while girls who enter puberty early experience negative outcomes. Boys who have earlier pubertal timing generally have more positive outcomes in adulthood but more negative outcomes in adolescence, while the reverse is true for later pubertal timing.[20]
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Outcomes have generally indicated that early onset of puberty in girls can be psychologically damaging. The main reason for this detrimental effect is the issue of body image. As they physically develop, gaining weight in several areas of the body, early-maturing girls usually look larger than girls who have not yet entered puberty. A result of the social pressure to be thin, the early-maturing girls develop a negative view of their body image. In addition, people may tease the girls about their visible breasts, forcing the early-maturing girl to hide her breasts by dressing differently. Embarrassment about a more developed body may also result in the refusal to undress for gym. These experiences lead to lower self-esteem, more depression and poorer body image in these early-maturing girls.[20]
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Furthermore, as physical and emotional differences set them apart from people in their same age group, early-maturing girls develop relationships with older people. For instance, some early-maturing girls have older boyfriends, "attracted to the girls' womanly physique and girlish innocence."[20] While having an older boyfriend might improve popularity among peers, it also increases the risk of alcohol and drug use, increased sexual relations (often unprotected), eating disorders and bullying.[20]
|
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+
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Generally, later onset of puberty in girls produces positive outcomes. They exhibit positive behaviors in adolescence that continue to adulthood.[20]
|
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In the past, early onset of puberty in boys has been associated with positive outcomes, such as leadership in high school and success in adulthood.[21] However, recent studies have revealed that the risks and problems of early maturation in males might outweigh the benefits.[20]
|
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Early-maturing boys develop "more aggressive, law-breaking, and alcohol abusing" behaviors, which result in anger towards parents and trouble in school and with the police. Early puberty also correlates with increased sexual activity and a higher instance of teenage pregnancy, both of which can lead to depression and other psychosocial issues.[20] However, early puberty might also result in positive outcomes, such as popularity among peers, higher self-esteem and confidence, as a result of physical developments, such as taller height, developed muscles, muscular male breast and better athletic ability.
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On the other hand, late-maturing boys develop lower self-esteem and confidence and generally have lower popularity among peers, due to their less-developed physiques. Also, they experience problems with anxiety and depression and are more likely to be afraid of sex than other boys.[20]
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In boys, puberty begins with the enlargement of the testicles and scrotum. The penis also increases in size, and a boy develops pubic hair. A boy's testicles also begin making sperm. The release of semen, which contains sperm and other fluids, is called ejaculation.[22] During puberty, a boy's erect penis becomes capable of ejaculating semen and impregnating a female.[23][24] A boy's first ejaculation is an important milestone in his development.[25] On average, a boy's first ejaculation occurs at age 13.[4] Ejaculation sometimes occurs during sleep; this phenomenon is known as a nocturnal emission.[22]
|
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In boys, testicular enlargement is the first physical manifestation of puberty (and is termed gonadarche).[26] Testes in prepubertal boys change little in size from about 1 year of age to the onset of puberty, averaging about 2–3 cm in length and about 1.5–2 cm in width. The size of the testicles is among the parameters of the Tanner scale for male genitals, from stage I which represents a volume of less than 1.5 ml, to stage V which represents a testicular volume of greater than 20 ml. Testicular size reaches maximal adult size about 6 years after the onset of puberty. While 18–20 cm3 is an average adult size, there is wide variation in testicular size in the normal population.[27] After the boy's testicles have enlarged and developed for about one year, the length and then the breadth of the shaft of the penis will increase and the glans penis and corpora cavernosa will also start to enlarge to adult proportions.[28]
|
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By the end of puberty, adult men have heavier bones and nearly twice as much skeletal muscle. Some of the bone growth (e.g. shoulder width and jaw) is disproportionately greater, resulting in noticeably different male and female skeletal shapes. The average adult male has about 150% of the lean body mass of an average female, and about 50% of the body fat.
|
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+
This muscle develops mainly during the later stages of puberty, and muscle growth can continue even after boys are biologically adult. The peak of the so-called "strength spurt", the rate of muscle growth, is attained about one year after a male experiences his peak growth rate.
|
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Often, the fat pads of the male breast tissue and the male nipples will develop during puberty; sometimes, especially in one breast, this becomes more apparent and is termed gynecomastia. It is usually not a permanent phenomenon.
|
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+
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Erections during sleep or when waking up are medically known as nocturnal penile tumescence and colloquially referred to as morning wood.[29] The penis can regularly get erect during sleep and men or boys often wake up with an erection.[30] Once a boy reaches his teenage years, erections occur much more frequently due to puberty.[31] Erections can occur spontaneously at any time of day, and if clothed may cause a bulge or "hump". This can be disguised or hidden by wearing close-fitting underwear, a long shirt and baggier clothes.[32] Erections are common for male prepubescent children and infants, and can even occur before birth.[33] Spontaneous erections are also known as involuntary or unwanted erections and are normal. Such erections can be embarrassing if they happen in public, such as a classroom or living room.[34][35]
|
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During puberty, if not before, the tip and opening of a boy's foreskin becomes wider, progressively allowing for retraction down the shaft of the penis and behind the glans, which ultimately should be possible without pain or difficulty. The membrane that bonds the inner surface of the foreskin with the glans disintegrates and releases the foreskin to separate from the glans. The foreskin then gradually becomes retractable.[36]
|
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|
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Research by Øster (1968) found that with the onset and continuation of puberty, the proportion of boys able to pull back their foreskins increased. At ages 12–13, Øster found that only 60% of boys were able to retract their foreskins; this increased to 85% by ages 14–15, and 95% by 16–17. He also found that 1% of those unable to fully retract experienced phimosis at ages 14–17, the remainder were partially able to.[36] The findings were supported by further research by Kayaba et al (1996) on a sample of over 600 boys,[37] and Ishikawa and Kawakita (2004) found that by age 15, 77% of their sample of boys could retract their foreskins.[38] Beaugé (1997) reports that boys may assist the development of retractile foreskin by manual stretching.[39]
|
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+
|
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Once a boy is able to retract his foreskin, penile hygiene should become an important feature of his routine body care. Although the American Academy of Pediatrics states there is "little evidence to affirm the association between circumcision status and optimal penile hygiene",[40] various studies suggest that boys be educated about the role of hygiene, including retracting the foreskin while urinating and rinsing under it and around the glans at each bathing opportunity. Regular washing under the foreskin was found by Krueger and Osborn (1986) to reduce the risk of numerous penile disorders,[41] however Birley et al. (1993) report excessive washing with soap should be avoided because it dries the oils out of the tissues and can cause non-specific dermatitis.[42]
|
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|
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+
Pubic hair often appears on a boy shortly after the genitalia begin to grow. The pubic hairs are usually first visible at the dorsal (abdominal) base of the penis. The first few hairs are described as stage 2. Stage 3 is usually reached within another 6–12 months, when the hairs are too many to count. By stage 4, the pubic hairs densely fill the "pubic triangle." Stage 5 refers to the spread of pubic hair to the thighs and upward towards the navel as part of the developing abdominal hair.
|
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|
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In the months and years following the appearance of pubic hair, other areas of skin that respond to androgens may develop androgenic hair. The usual sequence is: underarm (axillary) hair, perianal hair, upper lip hair, sideburn (preauricular) hair, periareolar hair, and the beard area.[43] As with most human biological processes, this specific order may vary among some individuals. Arm, leg, chest, abdominal, and back hair become heavier more gradually. There is a large range in amount of body hair among adult men, and significant differences in timing and quantity of hair growth among different racial groups. Facial hair is often present in late adolescence, but may not appear until significantly later.[44][45] Facial hair will continue to get coarser, darker and thicker for another 2–4 years after puberty.[44] Some men do not develop full facial hair for up to 10 years after the completion of puberty.[44] Chest hair may appear during puberty or years after, though not all men develop it.
|
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|
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+
Under the influence of androgens, the voice box, or larynx, grows in both sexes. This growth is far more prominent in boys, causing the male voice to drop and deepen, sometimes abruptly but rarely "overnight," about one octave, because the longer and thicker vocal folds have a lower fundamental frequency. Before puberty, the larynx of boys and girls is about equally small.[46] Occasionally, voice change is accompanied by unsteadiness of vocalization in the early stages of untrained voices. Most of the voice change happens during stage 3–4 of male puberty around the time of peak growth. Adult pitch is attained at an average age of 15 years, although the voice may not fully settle until early twenties. It usually precedes the development of significant facial hair by several months to years.
|
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|
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The first physical sign of puberty in girls is usually a firm, tender lump under the center of the areola of one or both breasts, occurring on average at about 10.5 years of age.[47] This is referred to as thelarche. By the widely used Tanner staging of puberty, this is stage 2 of breast development (stage 1 is a flat, prepubertal breast). Within six to 12 months, the swelling has clearly begun in both sides, softened, and can be felt and seen extending beyond the edges of the areolae. This is stage 3 of breast development. By another 12 months (stage 4), the breasts are approaching mature size and shape, with areolae and nipples forming a secondary mound. In most young women, this mound disappears into the contour of the mature breast (stage 5), although there is so much variation in sizes and shapes of adult breasts that stages 4 and 5 are not always separately identifiable.[48]
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|
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Pubic hair is often the second noticeable change in puberty, usually within a few months of thelarche.[49] It is referred to as pubarche. The pubic hairs are usually visible first along the labia. The first few hairs are described as Tanner stage 2.[48] Stage 3 is usually reached within another 6–12 months, when the hairs are too numerous to count and appear on the pubic mound as well. By stage 4, the pubic hairs densely fill the "pubic triangle." Stage 5 refers to spread of pubic hair to the thighs and sometimes as abdominal hair upward towards the navel. In about 15% of girls, the earliest pubic hair appears before breast development begins.[49]
|
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|
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Perineal skin keratinizes due to effect of estrogen increasing its resistance to infection. The mucosal surface of the vagina also changes in response to increasing levels of estrogen, becoming thicker and duller pink in color (in contrast to the brighter red of the prepubertal vaginal mucosa).[50] Mucosa changes into a multilayered structure with superficial layer of squamous cells. Estrogen increase glycogen content in vaginal epithelium, which in future plays important part in maintaining vaginal pH. Whitish secretions (physiologic leukorrhea) are a normal effect of estrogen as well.[47] In the two years following thelarche, the uterus, ovaries, and the follicles in the ovaries increase in size.[51] The ovaries usually contain small follicular cysts visible by ultrasound.[52][53] Before puberty, uterine body to cervix ratio is 1:1; which increases to 2:1 or 3:1 after completion of pubertal period.
|
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|
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The first menstrual bleeding is referred to as menarche, and typically occurs about two years after thelarche.[49] The average age of menarche is 12.5 in the United States.[54] Most American girls experience their first period at 11, 12 or 13, but some experience it earlier than their 11th birthday and others after their 14th birthday. In fact, anytime between 8 and 16 is normal. In Canada, the average age of menarche is 12.72,[55] and in the United Kingdom it is 12.9.[56] The time between menstrual periods (menses) is not always regular in the first two years after menarche.[57] Ovulation is necessary for fertility, but may or may not accompany the earliest menses.[58] In postmenarchal girls, about 80% of the cycles were anovulatory in the first year after menarche, 50% in the third year and 10% in the sixth year.[57] Initiation of ovulation after menarche is not inevitable. A high proportion of girls with continued irregularity in the menstrual cycle several years from menarche will continue to have prolonged irregularity and anovulation, and are at higher risk for reduced fertility.[59]
|
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|
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During this period, also in response to rising levels of estrogen, the lower half of the pelvis and thus hips widen (providing a larger birth canal).[48][60] Fat tissue increases to a greater percentage of the body composition than in males, especially in the typical female distribution of breasts, hips, buttocks, thighs, upper arms, and pubis. Progressive differences in fat distribution as well as sex differences in local skeletal growth contribute to the typical female body shape by the end of puberty. On average, at 10 years, girls have 6% more body fat than boys.[61]
|
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|
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Rising levels of androgens can change the fatty acid composition of perspiration, resulting in a more "adult" body odor. This often precedes thelarche and pubarche by one or more years. Another androgen effect is increased secretion of oil (sebum) from the skin. This change increases the susceptibility to acne, a skin condition that is characteristic of puberty.[62] Acne varies greatly in its severity.[62]
|
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|
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In girls, estradiol (the primary female sex hormone) causes thickening of lips and oral mucosa as well as further development of the vulva. In the vulva and vagina, estradiol causes thickening (stratification) of the skin and the growth of both the myoepithelial layer and the smooth muscle of the vagina. Typically estradiol will also cause pronounced growth of the labia minora and to a lesser degree of the labia majora.
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Estradiol is also responsible for the increased production of pheomelanin, resulting in the characteristic red color of the lips, labia minora and sometimes labia majora. Estradiol together with other ovarian steroids also cause the darker coloration of the areola.
|
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|
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Testosterone will cause an enlargement of the clitoris and possibly has important effects on the growth and maturation of the vestibular bulbs, corpus cavernosum of the clitoris and urethral sponge.[63]
|
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|
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Changes of the vulva initiated by estradiol as well as its direct effects also appear to influence the functioning of the lower urinary tract.[64][65]
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Hair growth develops under the arms, starting out sparse before thickening and darkening over time.[66]
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In a general sense, the conclusion of puberty is reproductive maturity. Criteria for defining the conclusion may differ for different purposes: attainment of the ability to reproduce, achievement of maximal adult height, maximal gonadal size, or adult sex hormone levels. Maximal adult height is achieved at an average age of 15 years for an average girl and 18 years for an average boy. Potential fertility (sometimes termed nubility) usually precedes completion of growth by 1–2 years in girls and 3–4 years in boys. Stage 5 typically represents maximal gonadal growth and adult hormone levels.
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The definition of the onset of puberty may depend on perspective (e.g., hormonal versus physical) and purpose (establishing population normal standards, clinical care of early or late pubescent individuals, etc.). A common definition for the onset of puberty is physical changes to a person's body.[10] These physical changes are the first visible signs of neural, hormonal, and gonadal function changes.
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The age at which puberty begins varies between individuals; usually, puberty begins between 10 and 13 years of age. The age at which puberty begins is affected by both genetic factors and by environmental factors such as nutritional state and social circumstances.[67] An example of social circumstances is the Vandenbergh effect; a juvenile female who has significant interaction with adult males will enter puberty earlier than juvenile females who are not socially overexposed to adult males.[68]
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The average age at which puberty begins may be affected by race as well. For example, the average age of menarche in various populations surveyed has ranged from 12[54][55][56] to 18 years. The earliest average onset of puberty is for African-American girls and the latest average onset for high altitude subsistence populations in Asia. However, much of the higher age averages reflect nutritional limitations more than genetic differences and can change within a few generations with a substantial change in diet. The median age of menarche for a population may be an index of the proportion of undernourished girls in the population, and the width of the spread may reflect unevenness of wealth and food distribution in a population.
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Researchers have identified an earlier age of the onset of puberty. However, they have based their conclusions on a comparison of data from 1999 with data from 1969. In the earlier example, the sample population was based on a small sample of white girls (200, from Britain). The later study identified as puberty as occurring in 48% of African-American girls by age nine, and 12% of white girls by that age.[69]
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One possible cause of a delay in the onset of puberty past the age 14 in girls and 15 in boys is Kallmann syndrome, a form of hypogonadotropic hypogonadism (HH). Kallmann syndrome is also associated with a lack of sense of smell (anosmia). Kallmann syndrome and other forms of HH affect both men and women. It is caused by a failure in HPG axis at puberty which results in low or zero gonadotropin (LH and FSH) levels with the subsequent result of a failure to commence or complete puberty, secondary hypogonadism and infertility.[70][71]
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Comparison of two individual with vast difference in the age of onset of puberty :
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The average age at which the onset of puberty occurs has dropped significantly since the 1840s.[72][73][74] In every decade from 1840 to 1950 there was a drop of four months in the average age of menarche among Western European females. In Norway, girls born in 1840 had their menarche at an average age of 17 years. In France, the average in 1840 was 15.3 years. In England, the average in 1840 was 16.5 years. In Japan the decline happened later and was then more rapid: from 1945 to 1975 in Japan there was a drop of 11 months per decade.
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A 2006 study in Denmark found that puberty, as evidenced by breast development, started at an average age of 9 years and 10 months, a year earlier than when a similar study was done in 1991. Scientists believe the phenomenon could be linked to obesity or exposure to chemicals in the food chain, and is putting girls at greater long-term risk of breast cancer.[75]
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Various studies have found direct genetic effects to account for at least 46% of the variation of timing of puberty in well-nourished populations.[76][77][78][79] The genetic association of timing is strongest between mothers and daughters. The specific genes affecting timing are not yet known.[76] Among the candidates is an androgen receptor gene.[80]
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Researchers[81] have hypothesized that early puberty onset may be caused by certain hair care products containing estrogen or placenta, and by certain chemicals, namely phthalates, which are used in many cosmetics, toys, and plastic food containers.
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If genetic factors account for half of the variation of pubertal timing, environment factors are clearly important as well. One of the first observed environmental effects is that puberty occurs later in children raised at higher altitudes. The most important of the environmental influences is clearly nutrition, but a number of others have been identified, all which affect timing of female puberty and menarche more clearly than male puberty.
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There is theoretical concern, and animal evidence, that environmental hormones and chemicals may affect aspects of prenatal or postnatal sexual development in humans.[82] Large amounts of incompletely metabolized estrogens and progestogens from pharmaceutical products are excreted into the sewage systems of large cities, and are sometimes detectable in the environment. Sex steroids are sometimes used in cattle farming but have been banned in chicken meat production for 40 years. Although agricultural laws regulate use to minimize accidental human consumption, the rules are largely self-enforced in the United States. Significant exposure of a child to hormones or other substances that activate estrogen or androgen receptors could produce some or all of the changes of puberty.
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Harder to detect as an influence on puberty are the more diffusely distributed environmental chemicals like PCBs (polychlorinated biphenyl), which can bind and trigger estrogen receptors.
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More obvious degrees of partial puberty from direct exposure of young children to small but significant amounts of pharmaceutical sex steroids from exposure at home may be detected during medical evaluation for precocious puberty, but mild effects and the other potential exposures outlined above would not.
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Bisphenol A (BPA) is a chemical used to make plastics, and is frequently used to make baby bottles, water bottles, sports equipment, medical devices, and as a coating in food and beverage cans. Scientists are concerned about BPA's behavioral effects on fetuses, infants, and children at current exposure levels because it can affect the prostate gland, mammary gland, and lead to early puberty in girls. BPA mimics and interferes with the action of estrogen—an important reproduction and development regulator. It leaches out of plastic into liquids and foods, and the Centers for Disease Control and Prevention (CDC) found measurable amounts of BPA in the bodies of more than 90 percent of the U.S. population studied. The highest estimated daily intakes of BPA occur in infants and children. Many plastic baby bottles contain BPA, and BPA is more likely to leach out of plastic when its temperature is increased, as when one warms a baby bottle or warms up food in the microwave.[83]
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Nutritional factors are the strongest and most obvious environmental factors affecting timing of puberty.[76] Girls are especially sensitive to nutritional regulation because they must contribute all of the nutritional support to a growing fetus. Surplus calories (beyond growth and activity requirements) are reflected in the amount of body fat, which signals to the brain the availability of resources for initiation of puberty and fertility.
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Much evidence suggests that for most of the last few centuries, nutritional differences accounted for majority of variation of pubertal timing in different populations, and even among social classes in the same population. Recent worldwide increased consumption of animal protein, other changes in nutrition, and increases in childhood obesity have resulted in falling ages of puberty, mainly in those populations with the higher previous ages. In many populations the amount of variation attributable to nutrition is shrinking.
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Although available dietary energy (simple calories) is the most important dietary influence on timing of puberty, quality of the diet plays a role as well. Lower protein intakes and higher dietary fiber intakes, as occur with typical vegetarian diets, are associated with later onset and slower progression of female puberty.
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Scientific researchers have linked early obesity with an earlier onset of puberty in girls. They have cited obesity as a cause of breast development before nine years and menarche before twelve years.[84] Early puberty in girls can be a harbinger of later health problems.[85]
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The average level of daily physical activity has also been shown to affect timing of puberty, especially in females. A high level of exercise, whether for athletic or body image purposes, or for daily subsistence, reduces energy calories available for reproduction and slows puberty. The exercise effect is often amplified by a lower body fat mass and cholesterol.
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Chronic diseases can delay puberty in both boys and girls. Those that involve chronic inflammation or interfere with nutrition have the strongest effect. In the western world, inflammatory bowel disease and tuberculosis have been notorious for such an effect in the last century, while in areas of the underdeveloped world, chronic parasite infections are widespread.
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Mental illnesses occur in puberty. The brain undergoes significant development by hormones which can contribute to mood disorders such as major depressive disorder, bipolar disorder, dysthymia and schizophrenia. Girls aged between 15 and 19 make up 40% of anorexia nervosa cases.[86]
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Some of the least understood environmental influences on timing of puberty are social and psychological. In comparison with the effects of genetics, nutrition, and general health, social influences are small, shifting timing by a few months rather than years. Mechanisms of these social effects are unknown, though a variety of physiological processes, including pheromones, have been suggested based on animal research.
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The most important part of a child's psychosocial environment is the family, and most of the social influence research has investigated features of family structure and function in relation to earlier or later female puberty. Most of the studies have reported that menarche may occur a few months earlier in girls in high-stress households, whose fathers are absent during their early childhood, who have a stepfather in the home, who are subjected to prolonged sexual abuse in childhood, or who are adopted from a developing country at a young age. Conversely, menarche may be slightly later when a girl grows up in a large family with a biological father present.
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More extreme degrees of environmental stress, such as wartime refugee status with threat to physical survival, have been found to be associated with delay of maturation, an effect that may be compounded by dietary inadequacy.
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Most of these reported social effects are small and our understanding is incomplete. Most of these "effects" are statistical associations revealed by epidemiologic surveys. Statistical associations are not necessarily causal, and a variety of covariables and alternative explanations can be imagined. Effects of such small size can never be confirmed or refuted for any individual child. Furthermore, interpretations of the data are politically controversial because of the ease with which this type of research can be used for political advocacy. Accusations of bias based on political agenda sometimes accompany scientific criticism.
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Another limitation of the social research is that nearly all of it has concerned girls, partly because female puberty requires greater physiologic resources and partly because it involves a unique event (menarche) that makes survey research into female puberty much simpler than male. More detail is provided in the menarche article.
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The sequence of events of pubertal development can occasionally vary. For example, in about 15% of boys and girls, pubarche (the first pubic hairs) can precede, respectively, gonadarche and thelarche by a few months. Rarely, menarche can occur before other signs of puberty in a few girls. These variations deserve medical evaluation because they can occasionally signal a disease.
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The endocrine reproductive system consists of the hypothalamus, the pituitary, the gonads, and the adrenal glands, with input and regulation from many other body systems. True puberty is often termed "central puberty" because it begins as a process of the central nervous system. A simple description of hormonal puberty is as follows:
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The onset of this neurohormonal process may precede the first visible body changes by 1–2 years.
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The arcuate nucleus of the hypothalamus is the driver of the reproductive system. It has neurons which generate and release pulses of GnRH into the portal venous system of the pituitary gland. The arcuate nucleus is affected and controlled by neuronal input from other areas of the brain and hormonal input from the gonads, adipose tissue and a variety of other systems.
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The pituitary gland responds to the pulsed GnRH signals by releasing LH and FSH into the blood of the general circulation, also in a pulsatile pattern.
|
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The gonads (testes and ovaries) respond to rising levels of LH and FSH by producing the steroid sex hormones, testosterone and estrogen.
|
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The adrenal glands are a second source for steroid hormones. Adrenal maturation, termed adrenarche, typically precedes gonadarche in mid-childhood.
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The endocrine reproductive system becomes functional by the end of the first trimester of fetal life. The testes and ovaries become briefly inactive around the time of birth but resume hormonal activity until several months after birth, when incompletely understood mechanisms in the brain begin to suppress the activity of the arcuate nucleus. This has been referred to as maturation of the prepubertal "gonadostat," which becomes sensitive to negative feedback by sex steroids. The period of hormonal activity until several months after birth, followed by suppression of activity, may correspond to the period of infant sexuality, followed by a latency stage, which Sigmund Freud described.[88]
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Gonadotropin and sex steroid levels fall to low levels (nearly undetectable by current clinical assays) for approximately another 8 to 10 years of childhood. Evidence is accumulating that the reproductive system is not totally inactive during the childhood years. Subtle increases in gonadotropin pulses occur, and ovarian follicles surrounding germ cells (future eggs) double in number.
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Normal puberty is initiated in the hypothalamus, with de-inhibition of the pulse generator in the arcuate nucleus. This inhibition of the arcuate nucleus is an ongoing active suppression by other areas of the brain. The signal and mechanism releasing the arcuate nucleus from inhibition have been the subject of investigation for decades and remain incompletely understood. Leptin levels rise throughout childhood and play a part in allowing the arcuate nucleus to resume operation. If the childhood inhibition of the arcuate nucleus is interrupted prematurely by injury to the brain, it may resume pulsatile gonadotropin release and puberty will begin at an early age.
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Neurons of the arcuate nucleus secrete gonadotropin releasing hormone (GnRH) into the blood of the pituitary portal system. An American physiologist, Ernst Knobil, found that the GnRH signals from the hypothalamus induce pulsed secretion of LH (and to a lesser degree, FSH) at roughly 1-2 hour intervals. The LH pulses are the consequence of pulsatile GnRH secretion by the arcuate nucleus that, in turn, is the result of an oscillator or signal generator in the central nervous system ("GnRH pulse generator").[89] In the years preceding physical puberty, Robert M. Boyar discovered that the gonadotropin pulses occur only during sleep, but as puberty progresses they can be detected during the day.[90] By the end of puberty, there is little day-night difference in the amplitude and frequency of gonadotropin pulses.
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Some investigators have attributed the onset of puberty to a resonance of oscillators in the brain.[91][92][93] By this mechanism, the gonadotropin pulses that occur primarily at night just before puberty represent beats.[94][95][96]
|
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An array of "autoamplification processes" increases the production of all of the pubertal hormones of the hypothalamus, pituitary, and gonads[citation needed].
|
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Regulation of adrenarche and its relationship to maturation of the hypothalamic-gonadal axis is not fully understood, and some evidence suggests it is a parallel but largely independent process coincident with or even preceding central puberty. Rising levels of adrenal androgens (termed adrenarche) can usually be detected between 6 and 11 years of age, even before the increasing gonadotropin pulses of hypothalamic puberty. Adrenal androgens contribute to the development of pubic hair (pubarche), adult body odor, and other androgenic changes in both sexes. The primary clinical significance of the distinction between adrenarche and gonadarche is that pubic hair and body odor changes by themselves do not prove that central puberty is underway for an individual child.
|
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|
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Early stages of male hypothalamic maturation seem to be very similar to the early stages of female puberty, though occurring about 1–2 years later.
|
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LH stimulates the Leydig cells of the testes to make testosterone and blood levels begin to rise. For much of puberty, nighttime levels of testosterone are higher than daytime. Regularity of frequency and amplitude of gonadotropin pulses seems to be less necessary for progression of male than female puberty.
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However, a significant portion of testosterone in adolescent boys is converted to estradiol. Estradiol mediates the growth spurt, bone maturation, and epiphyseal closure in boys just as in girls. Estradiol also induces at least modest development of breast tissue (gynecomastia) in a large proportion of boys. Boys who develop mild gynecomastia or even developing swellings under nipples during puberty are told the effects are temporary in some male teenagers due to high levels of estradiol.
|
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|
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Another hormonal change in males takes place during the teenage years for most young men. At this point in a male's life the testosterone levels slowly rise, and most of the effects are mediated through the androgen receptors by way of conversion dihydrotestosterone in target organs (especially that of the bowels).
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As the amplitude of LH pulses increases, the theca cells of the ovaries begin to produce testosterone and smaller amounts of progesterone. Much of the testosterone moves into nearby cells called granulosa cells. Smaller increases of FSH induce an increase in the aromatase activity of these granulosa cells, which converts most of the testosterone to estradiol for secretion into the circulation. The remaining testosterone, together with adrenal androgens is responsible for the typical androgenic changes of female puberty: pubic hair, other androgenic hair as outlined above, body odor, acne. The bioactivity of testosterone is to a large degree limited by SHBG which in turn is mainly controlled by estradiol and prolactin levels (estradiol stimulates, prolactin decreases SHBG synthesis).
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Rising levels of estradiol produce the characteristic estrogenic body changes of female puberty: growth spurt, acceleration of bone maturation and closure, breast growth, increased fat composition, growth of the uterus, increased thickness of the endometrium and the vaginal mucosa, and widening of the lower pelvis.
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As the estradiol levels gradually rise and the other autoamplification processes occur, a point of maturation is reached when the feedback sensitivity of the hypothalamic "gonadostat" becomes positive. This attainment of positive feedback is the hallmark of female sexual maturity, as it allows the mid cycle LH surge necessary for ovulation.
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Growth hormone levels rise steadily throughout puberty. IGF1 levels rise and then decline as puberty ends. Growth finishes and adult height is attained as the estradiol levels complete closure of the epiphyses.
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1 |
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Advertising is a marketing communication that employs an openly sponsored, non-personal message to promote or sell a product, service or idea.[1]:465 Sponsors of advertising are typically businesses wishing to promote their products or services. Advertising is differentiated from public relations in that an advertiser pays for and has control over the message. It differs from personal selling in that the message is non-personal, i.e., not directed to a particular individual.[1]:661,672
|
6 |
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Advertising is communicated through various mass media,[2] including traditional media such as newspapers, magazines, television, radio, outdoor advertising or direct mail; and new media such as search results, blogs, social media, websites or text messages. The actual presentation of the message in a medium is referred to as an advertisement, or "ad" or advert for short.
|
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Commercial ads often seek to generate increased consumption of their products or services through "branding", which associates a product name or image with certain qualities in the minds of consumers. On the other hand, ads that intend to elicit an immediate sale are known as direct-response advertising. Non-commercial entities that advertise more than consumer products or services include political parties, interest groups, religious organizations and governmental agencies. Non-profit organizations may use free modes of persuasion, such as a public service announcement. Advertising may also help to reassure employees or shareholders that a company is viable or successful.
|
9 |
+
|
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+
Modern advertising originated with the techniques introduced with tobacco advertising in the 1920s, most significantly with the campaigns of Edward Bernays, considered the founder of modern, "Madison Avenue" advertising.[3][4]
|
11 |
+
|
12 |
+
Worldwide spending on advertising in 2015 amounted to an estimated US$529.43 billion.[5] Advertising's projected distribution for 2017 was 40.4% on TV, 33.3% on digital, 9% on newspapers, 6.9% on magazines, 5.8% on outdoor and 4.3% on radio.[6] Internationally, the largest ("Big Five") advertising-agency groups are Dentsu, Interpublic, Omnicom, Publicis, and WPP.[7]
|
13 |
+
|
14 |
+
In Latin, advertere means "to turn towards".[8]
|
15 |
+
|
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+
Egyptians used papyrus to make sales messages and wall posters.[9] Commercial messages and political campaign displays have been found in the ruins of Pompeii and ancient Arabia. Lost and found advertising on papyrus was common in ancient Greece and ancient Rome. Wall or rock painting for commercial advertising is another manifestation of an ancient advertising form, which is present to this day in many parts of Asia, Africa, and South America. The tradition of wall painting can be traced back to Indian rock art paintings that date back to 4000 BC.[10]
|
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In ancient China, the earliest advertising known was oral, as recorded in the Classic of Poetry (11th to 7th centuries BC) of bamboo flutes played to sell confectionery. Advertisement usually takes in the form of calligraphic signboards and inked papers. A copper printing plate dated back to the Song dynasty used to print posters in the form of a square sheet of paper with a rabbit logo with "Jinan Liu's Fine Needle Shop" and "We buy high-quality steel rods and make fine-quality needles, to be ready for use at home in no time" written above and below[11] is considered the world's earliest identified printed advertising medium.[12]
|
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In Europe, as the towns and cities of the Middle Ages began to grow, and the general population was unable to read, instead of signs that read "cobbler", "miller", "tailor", or "blacksmith", images associated with their trade would be used such as a boot, a suit, a hat, a clock, a diamond, a horseshoe, a candle or even a bag of flour. Fruits and vegetables were sold in the city square from the backs of carts and wagons and their proprietors used street callers (town criers) to announce their whereabouts. The first compilation of such advertisements was gathered in "Les Crieries de Paris", a thirteenth-century poem by Guillaume de la Villeneuve.[13]
|
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+
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+
In the 18th century advertisements started to appear in weekly newspapers in England. These early print advertisements were used mainly to promote books and newspapers, which became increasingly affordable with advances in the printing press; and medicines, which were increasingly sought after. However, false advertising and so-called "quack" advertisements became a problem, which ushered in the regulation of advertising content.
|
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+
Thomas J. Barratt of London has been called "the father of modern advertising".[14][15][16] Working for the Pears Soap company, Barratt created an effective advertising campaign for the company products, which involved the use of targeted slogans, images and phrases. One of his slogans, "Good morning. Have you used Pears' soap?" was famous in its day and into the 20th century.[17][18]
|
25 |
+
|
26 |
+
Barratt introduced many of the crucial ideas that lie behind successful advertising and these were widely circulated in his day. He constantly stressed the importance of a strong and exclusive brand image for Pears and of emphasizing the product's availability through saturation campaigns. He also understood the importance of constantly reevaluating the market for changing tastes and mores, stating in 1907 that "tastes change, fashions change, and the advertiser has to change with them. An idea that was effective a generation ago would fall flat, stale, and unprofitable if presented to the public today. Not that the idea of today is always better than the older idea, but it is different – it hits the present taste."[15]
|
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|
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+
As the economy expanded across the world during the 19th century, advertising grew alongside. In the United States, the success of this advertising format eventually led to the growth of mail-order advertising.
|
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+
|
30 |
+
In June 1836, French newspaper La Presse was the first to include paid advertising in its pages, allowing it to lower its price, extend its readership and increase its profitability and the formula was soon copied by all titles. Around 1840, Volney B. Palmer established the roots of the modern day advertising agency in Philadelphia. In 1842 Palmer bought large amounts of space in various newspapers at a discounted rate then resold the space at higher rates to advertisers. The actual ad – the copy, layout, and artwork – was still prepared by the company wishing to advertise; in effect, Palmer was a space broker. The situation changed when the first full-service advertising agency of N.W. Ayer & Son was founded in 1869 in Philadelphia. Ayer & Son offered to plan, create, and execute complete advertising campaigns for its customers. By 1900 the advertising agency had become the focal point of creative planning, and advertising was firmly established as a profession.
|
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+
[19] Around the same time, in France, Charles-Louis Havas extended the services of his news agency, Havas to include advertisement brokerage, making it the first French group to organize. At first, agencies were brokers for advertisement space in newspapers.[19]
|
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+
|
33 |
+
Advertising increased dramatically in the United States as industrialization expanded the supply of manufactured products. In 1919 it was 2.5 percent of gross domestic product (GDP) in the US, and it averaged 2.2 percent of GDP between then and at least 2007, though it may have declined dramatically since the Great Recession.
|
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+
|
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+
Industry could not benefit from its increased productivity without a substantial increase in consumer spending. This contributed to the development of mass marketing designed to influence the population's economic behavior on a larger scale.[20] In the 1910s and 1920s, advertisers in the U.S. adopted the doctrine that human instincts could be targeted and harnessed – "sublimated" into the desire to purchase commodities.[21] Edward Bernays, a nephew of Sigmund Freud, became associated with the method and is sometimes called the founder of modern advertising and public relations.[22] Bernays claimed that:
|
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+
|
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+
"[The] general principle, that men are very largely actuated by motives which they conceal from themselves, is as true of mass as of individual psychology. It is evident that the successful propagandist must understand the true motives and not be content to accept the reasons which men give for what they do."[23]
|
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+
|
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+
In other words, selling products by appealing to the rational minds of customers (the main method used prior to Bernays) was much less effective than selling products based on the unconscious desires that Bernays felt were the true motivators of human action. "Sex sells" became a controversial issue, with techniques for titillating and enlarging the audience posing a challenge to conventional morality.[24][25]
|
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+
|
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+
In the 1920s, under Secretary of Commerce Herbert Hoover, the American government promoted advertising. Hoover himself delivered an address to the Associated Advertising Clubs of the World in 1925 called 'Advertising Is a Vital Force in Our National Life."[26] In October 1929, the head of the U.S. Bureau of Foreign and Domestic Commerce, Julius Klein, stated "Advertising is the key to world prosperity."[27] This was part of the "unparalleled" collaboration between business and government in the 1920s, according to a 1933 European economic journal.[28]
|
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+
|
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+
The tobacco companies became major advertisers in order to sell packaged cigarettes.[29] The tobacco companies pioneered the new advertising techniques when they hired Bernays to create positive associations with tobacco smoking.[3][4]
|
44 |
+
|
45 |
+
Advertising was also used as a vehicle for cultural assimilation, encouraging workers to exchange their traditional habits and community structure in favor of a shared "modern" lifestyle.[30] An important tool for influencing immigrant workers was the American Association of Foreign Language Newspapers (AAFLN). The AAFLN was primarily an advertising agency but also gained heavily centralized control over much of the immigrant press.[31][32]
|
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+
|
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+
At the turn of the 20th century, advertising was one of the few career choices for women. Since women were responsible for most household purchasing done, advertisers and agencies recognized the value of women's insight during the creative process. In fact, the first American advertising to use a sexual sell was created by a woman – for a soap product. Although tame by today's standards, the advertisement featured a couple with the message "A skin you love to touch".[33]
|
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+
|
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+
In the 1920s psychologists Walter D. Scott and John B. Watson contributed applied psychological theory to the field of advertising. Scott said, "Man has been called the reasoning animal but he could with greater truthfulness be called the creature of suggestion. He is reasonable, but he is to a greater extent suggestible".[34] He demonstrated this through his advertising technique of a direct command to the consumer.
|
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+
|
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+
In the early 1920s, the first radio stations were established by radio equipment manufacturers, followed by non-profit organizations such as schools, clubs and civic groups who also set up their own stations.[35] Retailer and consumer goods manufacturers quickly recognised radio's potential to reach consumers in their home and soon adopted advertising techniques that would allow their messages to stand out; slogans, mascots, and jingles began to appear on radio in the 1920s and early television in the 1930s.[36]
|
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|
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+
The rise of mass media communications allowed manufacturers of branded goods to bypass retailers by advertising directly to consumers. This was a major paradigm shift which forced manufacturers to focus on the brand and stimulated the need for superior insights into consumer purchasing, consumption and usage behaviour; their needs, wants and aspirations.[37] The earliest radio drama series were sponsored by soap manufacturers and the genre became known as a soap opera.[38] Before long, radio station owners realized they could increase advertising revenue by selling 'air-time' in small time allocations which could be sold to multiple businesses. By the 1930s, these advertising spots, as the packets of time became known, were being sold by the station's geographical sales representatives, ushering in an era of national radio advertising.[39]
|
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+
By the 1940s, manufacturers began to recognize the way in which consumers were developing personal relationships with their brands in a social/psychological/anthropological sense.[40] Advertisers began to use motivational research and consumer research to gather insights into consumer purchasing. Strong branded campaigns for Chrysler and Exxon/Esso, using insights drawn research methods from psychology and cultural anthropology, led to some of the most enduring campaigns of the 20th-century.[41]
|
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|
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+
In the early 1950s, the DuMont Television Network began the modern practice of selling advertisement time to multiple sponsors. Previously, DuMont had trouble finding sponsors for many of their programs and compensated by selling smaller blocks of advertising time to several businesses. This eventually became the standard for the commercial television industry in the United States. However, it was still a common practice to have single sponsor shows, such as The United States Steel Hour. In some instances the sponsors exercised great control over the content of the show – up to and including having one's advertising agency actually writing the show.[citation needed] The single sponsor model is much less prevalent now, a notable exception being the Hallmark Hall of Fame.[citation needed]
|
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|
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+
The late 1980s and early 1990s saw the introduction of cable television and particularly MTV. Pioneering the concept of the music video, MTV ushered in a new type of advertising: the consumer tunes in for the advertising message, rather than it being a by-product or afterthought. As cable and satellite television became increasingly prevalent, specialty channels emerged, including channels entirely devoted to advertising, such as QVC, Home Shopping Network, and ShopTV Canada.[42]
|
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|
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+
With the advent of the ad server, online advertising grew, contributing to the "dot-com" boom of the 1990s.[43] Entire corporations operated solely on advertising revenue, offering everything from coupons to free Internet access. At the turn of the 21st century, some websites, including the search engine Google, changed online advertising by personalizing ads based on web browsing behavior. This has led to other similar efforts and an increase in interactive advertising.[44]
|
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The share of advertising spending relative to GDP has changed little across large changes in media since 1925. In 1925, the main advertising media in America were newspapers, magazines, signs on streetcars, and outdoor posters. Advertising spending as a share of GDP was about 2.9 percent. By 1998, television and radio had become major advertising media; by 2017, the balance between broadcast and online advertising had shifted, with online spending exceeding broadcast.[45] Nonetheless, advertising spending as a share of GDP was slightly lower – about 2.4 percent.[46]
|
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Guerrilla marketing involves unusual approaches such as staged encounters in public places, giveaways of products such as cars that are covered with brand messages, and interactive advertising where the viewer can respond to become part of the advertising message. This type of advertising is unpredictable, which causes consumers to buy the product or idea.[47] This reflects an increasing trend of interactive and "embedded" ads, such as via product placement, having consumers vote through text messages, and various campaigns utilizing social network services such as Facebook or Twitter.[48]
|
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+
The advertising business model has also been adapted in recent years.[when?][clarification needed] In media for equity, advertising is not sold, but provided to start-up companies in return for equity. If the company grows and is sold, the media companies receive cash for their shares.
|
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Domain name registrants (usually those who register and renew domains as an investment) sometimes "park" their domains and allow advertising companies to place ads on their sites in return for per-click payments. These ads are typically driven by pay per click search engines like Google or Yahoo, but ads can sometimes be placed directly on targeted domain names through a domain lease or by making contact with the registrant of a domain name that describes a product. Domain name registrants are generally easy to identify through WHOIS records that are publicly available at registrar websites.[49]
|
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Advertising may be categorized in a variety of ways, including by style, target audience, geographic scope, medium, or purpose.[2]:9–15 For example, in print advertising, classification by style can include display advertising (ads with design elements sold by size) vs. classified advertising (ads without design elements sold by the word or line). Advertising may be local, national or global. An ad campaign may be directed toward consumers or to businesses. The purpose of an ad may be to raise awareness (brand advertising), or to elicit an immediate sale (direct response advertising). The term above the line (ATL) is used for advertising involving mass media; more targeted forms of advertising and promotion are referred to as below the line (BTL).[50][51] The two terms date back to 1954 when Procter & Gamble began paying their advertising agencies differently from other promotional agencies.[52] In the 2010s, as advertising technology developed, a new term, through the line (TTL) began to come into use, referring to integrated advertising campaigns.[53][54]
|
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Virtually any medium can be used for advertising. Commercial advertising media can include wall paintings, billboards, street furniture components, printed flyers and rack cards, radio, cinema and television adverts, web banners, mobile telephone screens, shopping carts, web popups, skywriting, bus stop benches, human billboards and forehead advertising, magazines, newspapers, town criers, sides of buses, banners attached to or sides of airplanes ("logojets"), in-flight advertisements on seatback tray tables or overhead storage bins, taxicab doors, roof mounts and passenger screens, musical stage shows, subway platforms and trains, elastic bands on disposable diapers, doors of bathroom stalls, stickers on apples in supermarkets, shopping cart handles (grabertising), the opening section of streaming audio and video, posters, and the backs of event tickets and supermarket receipts. Any situation in which an "identified" sponsor pays to deliver their message through a medium is advertising.[55]
|
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|
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A new advertising approach is known as advanced advertising, which is data-driven advertising, using large quantities of data, precise measuring tools and precise targeting.[66] Advanced advertising also makes it easier for companies which sell ad-space to attribute customer purchases to the ads they display or broadcast.[67]
|
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|
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Increasingly, other media are overtaking many of the "traditional" media such as television, radio and newspaper because of a shift toward the usage of the Internet for news and music as well as devices like digital video recorders (DVRs) such as TiVo.[68]
|
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|
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Online advertising began with unsolicited bulk e-mail advertising known as "e-mail spam". Spam has been a problem for e-mail users since 1978.[69] As new online communication channels became available, advertising followed. The first banner ad appeared on the World Wide Web in 1994.[70] Prices of Web-based advertising space are dependent on the "relevance" of the surrounding web content and the traffic that the website receives.[citation needed]
|
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In online display advertising, display ads generate awareness quickly. Unlike search, which requires someone to be aware of a need, display advertising can drive awareness of something new and without previous knowledge. Display works well for direct response. Display is not only used for generating awareness, it's used for direct response campaigns that link to a landing page with a clear 'call to action'.[71]
|
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|
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As the mobile phone became a new mass medium in 1998 when the first paid downloadable content appeared on mobile phones in Finland,[72][citation needed] mobile advertising followed, also first launched in Finland in 2000.[citation needed] By 2007 the value of mobile advertising had reached $2 billion and providers such as Admob delivered billions of mobile ads.[citation needed]
|
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|
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More advanced mobile ads include banner ads, coupons, Multimedia Messaging Service picture and video messages, advergames and various engagement marketing campaigns. A particular feature driving mobile ads is the 2D barcode, which replaces the need to do any typing of web addresses, and uses the camera feature of modern phones to gain immediate access to web content. 83 percent of Japanese mobile phone users already are active users of 2D barcodes.[citation needed]
|
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|
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+
Some companies have proposed placing messages or corporate logos on the side of booster rockets and the International Space Station.[citation needed]
|
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+
|
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Unpaid advertising (also called "publicity advertising"), can include personal recommendations ("bring a friend", "sell it"), spreading buzz, or achieving the feat of equating a brand with a common noun (in the United States, "Xerox" = "photocopier", "Kleenex" = tissue, "Vaseline" = petroleum jelly, "Hoover" = vacuum cleaner, and "Band-Aid" = adhesive bandage). However, some companies[which?] oppose the use of their brand name to label an object. Equating a brand with a common noun also risks turning that brand into a generic trademark – turning it into a generic term which means that its legal protection as a trademark is lost.[73][disputed – discuss]
|
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From time to time, The CW Television Network airs short programming breaks called "Content Wraps", to advertise one company's product during an entire commercial break. The CW pioneered "content wraps" and some products featured were Herbal Essences, Crest, Guitar Hero II, CoverGirl, and Toyota.[74][75]
|
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A new promotion concept has appeared, "ARvertising", advertising on augmented reality technology.[76]
|
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+
|
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Controversy exists on the effectiveness of subliminal advertising (see mind control), and the pervasiveness of mass messages (propaganda).
|
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|
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With the Internet came many new advertising opportunities. Pop-up, Flash, banner, pop-under, advergaming, and email advertisements (all of which are often unwanted or spam in the case of email) are now commonplace. Particularly since the rise of "entertaining" advertising, some people may like an advertisement enough to wish to watch it later or show a friend.[citation needed] In general, the advertising community has not yet made this easy, although some have used the Internet to widely distribute their ads to anyone willing to see or hear them. In the last three quarters of 2009, mobile and Internet advertising grew by 18% and 9% respectively, while older media advertising saw declines: −10.1% (TV), −11.7% (radio), −14.8% (magazines) and −18.7% (newspapers).[citation needed] Between 2008 and 2014, U.S. newspapers lost more than half their print advertising revenue.[78]
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Another significant trend regarding future of advertising is the growing importance of the niche market using niche or targeted ads. Also brought about by the Internet and the theory of the long tail, advertisers will have an increasing ability to reach specific audiences. In the past, the most efficient way to deliver a message was to blanket the largest mass market audience possible.[citation needed] However, usage tracking, customer profiles and the growing popularity of niche content brought about by everything from blogs to social networking sites, provide advertisers with audiences that are smaller but much better defined,[citation needed] leading to ads that are more relevant to viewers and more effective for companies' marketing products. Among others, Comcast Spotlight is one such advertiser employing this method in their video on demand menus. These advertisements are targeted to a specific group and can be viewed by anyone wishing to find out more about a particular business or practice, from their home. This causes the viewer to become proactive and actually choose what advertisements they want to view.[79]
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Niche marketing could also be helped by bringing the issue of colour into advertisements. Different colours play major roles when it comes to marketing strategies, for example, seeing the blue can promote a sense of calmness and gives a sense of security which is why many social networks such as Facebook use blue in their logos.
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Google AdSense is an example of niche marketing. Google calculates the primary purpose of a website and adjusts ads accordingly; it uses keywords on the page (or even in emails) to find the general ideas of topics disused and places ads that will most likely be clicked on by viewers of the email account or website visitors.
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The concept of crowdsourcing has given way to the trend of user-generated advertisements. User-generated ads are created by people, as opposed to an advertising agency or the company themselves, often resulting from brand sponsored advertising competitions. For the 2007 Super Bowl, the Frito-Lays division of PepsiCo held the "Crash the Super Bowl" contest, allowing people to create their own Doritos commercials.[80] Chevrolet held a similar competition for their Tahoe line of SUVs.[80] Due to the success of the Doritos user-generated ads in the 2007 Super Bowl, Frito-Lays relaunched the competition for the 2009 and 2010 Super Bowl. The resulting ads were among the most-watched and most-liked Super Bowl ads. In fact, the winning ad that aired in the 2009 Super Bowl was ranked by the USA Today Super Bowl Ad Meter as the top ad for the year while the winning ads that aired in the 2010 Super Bowl were found by Nielsen's BuzzMetrics to be the "most buzzed-about".[81][82] Another example of companies using crowdsourcing successfully is the beverage company Jones Soda that encourages consumers to participate in the label design themselves.[83]
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This trend has given rise to several online platforms that host user-generated advertising competitions on behalf of a company. Founded in 2007, Zooppa has launched ad competitions for brands such as Google, Nike, Hershey's, General Mills, Microsoft, NBC Universal, Zinio, and Mini Cooper.[84] Crowdsourcing remains controversial, as the long-term impact on the advertising industry is still unclear.[85]
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Advertising has gone through five major stages of development: domestic, export, international, multi-national, and global. For global advertisers, there are four, potentially competing, business objectives that must be balanced when developing worldwide advertising: building a brand while speaking with one voice, developing economies of scale in the creative process, maximising local effectiveness of ads, and increasing the company's speed of implementation. Born from the evolutionary stages of global marketing are the three primary and fundamentally different approaches to the development of global advertising executions: exporting executions, producing local executions, and importing ideas that travel.[86]
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Advertising research is key to determining the success of an ad in any country or region. The ability to identify which elements and/or moments of an ad contribute to its success is how economies of scale are maximized. Once one knows what works in an ad, that idea or ideas can be imported by any other market. Market research measures, such as Flow of Attention, Flow of Emotion and branding moments provide insight into what is working in an ad in any country or region because the measures are based on the visual, not verbal, elements of the ad.[87]
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Foreign governments,[which?] particularly those that own marketable commercial products or services, often promote their interests and positions through the advertising of those goods because the target audience is not only largely unaware of the forum as a vehicle for foreign messaging but also willing to receive the message while in a mental state of absorbing information from advertisements during television commercial breaks, while reading a periodical, or while passing by billboards in public spaces. A prime example of this messaging technique is advertising campaigns to promote international travel. While advertising foreign destinations and services may stem from the typical goal of increasing revenue by drawing more tourism, some travel campaigns carry the additional or alternative intended purpose of promoting good sentiments or improving existing ones among the target audience towards a given nation or region. It is common for advertising promoting foreign countries to be produced and distributed by the tourism ministries of those countries, so these ads often carry political statements and/or depictions of the foreign government's desired international public perception. Additionally, a wide range of foreign airlines and travel-related services which advertise separately from the destinations, themselves, are owned by their respective governments; examples include, though are not limited to, the Emirates airline (Dubai), Singapore Airlines (Singapore), Qatar Airways (Qatar), China Airlines (Taiwan/Republic of China), and Air China (People's Republic of China). By depicting their destinations, airlines, and other services in a favorable and pleasant light, countries market themselves to populations abroad in a manner that could mitigate prior public impressions.
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In the realm of advertising agencies, continued industry diversification has seen observers note that "big global clients don't need big global agencies any more".[88] This is reflected by the growth of non-traditional agencies in various global markets, such as Canadian business TAXI and SMART in Australia and has been referred to as "a revolution in the ad world".[89]
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The ability to record shows on digital video recorders (such as TiVo) allow watchers to record the programs for later viewing, enabling them to fast forward through commercials. Additionally, as more seasons of pre-recorded box sets are offered for sale of television programs; fewer people watch the shows on TV. However, the fact that these sets are sold, means the company will receive additional profits from these sets.
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To counter this effect, a variety of strategies have been employed. Many advertisers have opted for product placement on TV shows like Survivor. Other strategies include integrating advertising with internet-connected program guidess (EPGs), advertising on companion devices (like smartphones and tablets) during the show, and creating mobile apps for TV programs. Additionally, some like brands have opted for social television sponsorship.[90]
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The emerging technology of drone displays has recently been used for advertising purposes.[91]
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In recent years there have been several media literacy initiatives, and more specifically concerning advertising, that seek to empower citizens in the face of media advertising campaigns.[92]
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Advertising education has become popular with bachelor, master and doctorate degrees becoming available in the emphasis.[citation needed] A surge in advertising interest is typically attributed to the strong relationship advertising plays in cultural and technological changes, such as the advance of online social networking.[citation needed] A unique model for teaching advertising is the student-run advertising agency, where advertising students create campaigns for real companies.[93] Organizations such as the American Advertising Federation establish companies with students to create these campaigns.[citation needed]
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Advertising is at the front of delivering the proper message to customers and prospective customers. The purpose of advertising is to inform the consumers about their product and convince customers that a company's services or products are the best, enhance the image of the company, point out and create a need for products or services, demonstrate new uses for established products, announce new products and programs, reinforce the salespeople's individual messages, draw customers to the business, and to hold existing customers.[94]
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Sales promotions are another way to advertise. Sales promotions are double purposed because they are used to gather information about what type of customers one draws in and where they are, and to jump start sales. Sales promotions include things like contests and games, sweepstakes, product giveaways, samples coupons, loyalty programs, and discounts. The ultimate goal of sales promotions is to stimulate potential customers to action.[95]
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While advertising can be seen as necessary for economic growth,[27] it is not without social costs. Unsolicited commercial e-mail and other forms of spam have become so prevalent as to have become a major nuisance to users of these services, as well as being a financial burden on internet service providers.[96] Advertising is increasingly invading public spaces, such as schools, which some critics argue is a form of child exploitation.[97] This increasing difficulty in limiting exposure to specific audiences can result in negative backlash for advertisers.[98] In tandem with these criticisms, the advertising industry has seen low approval rates in surveys and negative cultural portrayals.[99]
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One of the most controversial criticisms of advertisement in the present day is that of the predominance of advertising of foods high in sugar, fat, and salt specifically to children. Critics claim that food advertisements targeting children are exploitive and are not sufficiently balanced with proper nutritional education to help children understand the consequences of their food choices. Additionally, children may not understand that they are being sold something, and are therefore more impressionable.[100] Michelle Obama has criticized large food companies for advertising unhealthy foods largely towards children and has requested that food companies either limit their advertising to children or advertise foods that are more in line with dietary guidelines.[101] The other criticisms include the change that are brought by those advertisements on the society and also the deceiving ads that are aired and published by the corporations. Cosmetic and health industry are the ones which exploited the highest and created reasons of concern.[102]
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There have been increasing efforts to protect the public interest by regulating the content and the influence of advertising. Some examples include restrictions for advertising alcohol, tobacco or gambling imposed in many countries, as well as the bans around advertising to children, which exist in parts of Europe. Advertising regulation focuses heavily on the veracity of the claims and as such, there are often tighter restrictions placed around advertisements for food and healthcare products.[103]
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The advertising industries within some countries rely less on laws and more on systems of self-regulation.[103][104][105] Advertisers and the media agree on a code of advertising standards that they attempt to uphold. The general aim of such codes is to ensure that any advertising is 'legal, decent, honest and truthful'. Some self-regulatory organizations are funded by the industry, but remain independent, with the intent of upholding the standards or codes like the Advertising Standards Authority in the UK.[106]
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In the UK, most forms of outdoor advertising such as the display of billboards is regulated by the UK Town and County Planning system. Currently, the display of an advertisement without consent from the Planning Authority is a criminal offense liable to a fine of £2,500 per offense.[107] In the US, many communities believe that many forms of outdoor advertising blight the public realm.[108] As long ago as the 1960s in the US, there were attempts to ban billboard advertising in the open countryside.[109] Cities such as São Paulo have introduced an outright ban[110] with London also having specific legislation to control unlawful displays.
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Some governments restrict the languages that can be used in advertisements, but advertisers may employ tricks to try avoiding them. In France for instance, advertisers sometimes print English words in bold and French translations in fine print to deal with Article 120 of the 1994 Toubon Law limiting the use of English.[111]
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The advertising of pricing information is another topic of concern for governments. In the United States for instance, it is common for businesses to only mention the existence and amount of applicable taxes at a later stage of a transaction.[112] In Canada and New Zealand, taxes can be listed as separate items, as long as they are quoted up-front.[113][114] In most other countries, the advertised price must include all applicable taxes, enabling customers to easily know how much it will cost them.[115][116][117]
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Various competing models of hierarchies of effects attempt to provide a theoretical underpinning to advertising practice.[clarification needed][118]
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The marketing mix was proposed by professor E. Jerome McCarthy in the 1960s.[122] It consists of four basic elements called the "four Ps". Product is the first P representing the actual product. Price represents the process of determining the value of a product. Place represents the variables of getting the product to the consumer such as distribution channels, market coverage and movement organization. The last P stands for Promotion which is the process of reaching the target market and convincing them to buy the product.
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In the 1990s, the concept of four Cs was introduced as a more customer-driven replacement of four P's.[123] There are two theories based on four Cs: Lauterborn's four Cs (consumer, cost, communication, convenience)
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[124] and Shimizu's four Cs (commodity, cost, communication, channel) in the 7Cs Compass Model (Co-marketing). Communications can include advertising, sales promotion, public relations, publicity, personal selling, corporate identity, internal communication, SNS, and MIS.[125][126][127][128]
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Advertising research is a specialized form of research that works to improve the effectiveness and efficiency of advertising. It entails numerous forms of research which employ different methodologies. Advertising research includes pre-testing (also known as copy testing) and post-testing of ads and/or campaigns.
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Pre-testing includes a wide range of qualitative and quantitative techniques, including: focus groups, in-depth target audience interviews (one-on-one interviews), small-scale quantitative studies and physiological measurement. The goal of these investigations is to better understand how different groups respond to various messages and visual prompts, thereby providing an assessment of how well the advertisement meets its communications goals.[129]
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Post-testing employs many of the same techniques as pre-testing, usually with a focus on understanding the change in awareness or attitude attributable to the advertisement. With the emergence of digital advertising technologies, many firms have begun to continuously post-test ads using real-time data. This may take the form of A/B split-testing or multivariate testing.
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Continuous ad tracking and the Communicus System are competing examples of post-testing advertising research types.[130]
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Meanings between consumers and marketers depict signs and symbols that are encoded in everyday objects.[131] Semiotics is the study of signs and how they are interpreted. Advertising has many hidden signs and meanings within brand names, logos, package designs, print advertisements, and television advertisements. Semiotics aims to study and interpret the message being conveyed in (for example) advertisements. Logos and advertisements can be interpreted at two levels – known as the surface level and the underlying level. The surface level uses signs creatively to create an image or personality for a product.[citation needed] These signs can be images, words, fonts, colors, or slogans. The underlying level is made up of hidden meanings. The combination of images, words, colors, and slogans must be interpreted by the audience or consumer.[132] The "key to advertising analysis" is the signifier and the signified. The signifier is the object and the signified is the mental concept.[133] A product has a signifier and a signified. The signifier is the color, brand name, logo design, and technology. The signified has two meanings known as denotative and connotative. The denotative meaning is the meaning of the product. A television's denotative meaning might be that it is high definition. The connotative meaning is the product's deep and hidden meaning. A connotative meaning of a television would be that it is top-of-the-line.[134]
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Apple's commercials[when?] used a black silhouette of a person that was the age of Apple's target market. They placed the silhouette in front of a blue screen so that the picture behind the silhouette could be constantly changing. However, the one thing that stays the same in these ads is that there is music in the background and the silhouette is listening to that music on a white iPod through white headphones. Through advertising, the white color on a set of earphones now signifies that the music device is an iPod. The white color signifies almost all of Apple's products.[135]
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The semiotics of gender plays a key influence on the way in which signs are interpreted. When considering gender roles in advertising, individuals are influenced by three categories. Certain characteristics of stimuli may enhance or decrease the elaboration of the message (if the product is perceived as feminine or masculine). Second, the characteristics of individuals can affect attention and elaboration of the message (traditional or non-traditional gender role orientation). Lastly, situational factors may be important to influence the elaboration of the message.[136]
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There are two types of marketing communication claims-objective and subjective.[137] Objective claims stem from the extent to which the claim associates the brand with a tangible product or service feature. For instance, a camera may have auto-focus features. Subjective claims convey emotional, subjective, impressions of intangible aspects of a product or service. They are non-physical features of a product or service that cannot be directly perceived, as they have no physical reality. For instance the brochure has a beautiful design.[138] Males tend to respond better to objective marketing-communications claims while females tend to respond better to subjective marketing communications claims.[139]
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Voiceovers are commonly used in advertising. Most voiceovers are done by men, with figures of up to 94% having been reported.[140] There have been more female voiceovers in recent years,[when?] but mainly for food, household products, and feminine-care products.[141]
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According to a 1977 study by David Statt, females process information comprehensively, while males process information through heuristic devices such as procedures, methods or strategies for solving problems, which could have an effect on how they interpret advertising.[142][need quotation to verify] According to this study, men prefer to have available and apparent cues to interpret the message, whereas females engage in more creative, associative, imagery-laced interpretation. Later research by a Danish team[143] found that advertising attempts to persuade men to improve their appearance or performance, whereas its approach to women aims at transformation toward an impossible ideal of female presentation. In Paul Suggett's article "The Objectification of Women in Advertising"[144] he discusses the negative impact that these women in advertisements, who are too perfect to be real, have on women, as well as men, in real life. Advertising's manipulation of women's aspiration to these ideal types as portrayed in film, in erotic art, in advertising, on stage, within music videos, and through other media exposures requires at least a conditioned rejection of female reality, and thereby takes on a highly ideological cast. Studies show that these expectations of women and young girls negatively impact their views about their bodies and appearances. These advertisements are directed towards men. Not everyone agrees: one critic viewed this monologic, gender-specific interpretation of advertising as excessively skewed and politicized.[145][need quotation to verify] There are some companies, however, like Dove and Aerie that are creating commercials to portray more natural women, with less post production manipulation, so more women and young girls are able to relate to them.[citation needed]
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More recent research by Martin (2003) reveals that males and females differ in how they react to advertising depending on their mood at the time of exposure to the ads, and on the affective tone of the advertising. When feeling sad, males prefer happy ads to boost their mood. In contrast, females prefer happy ads when they are feeling happy. The television programs in which ads are embedded influence a viewer's mood state.[146] Susan Wojcicki, author of the article "Ads that Empower Women don’t just Break Stereotypes—They’re also Effective"[147] discusses how advertising to women has changed since the first Barbie commercial where a little girl tells the doll that, she wants to be just like her. Little girls grow up watching advertisements of scantily clad women advertising things from trucks to burgers, and Wojcicki states that this shows girls that they are either arm candy or eye candy.
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Other approaches to revenue include donations, paid subscriptions and microtransactions. Websites and applications are "ad-free" when not using ads at all for revenue. For example, the online encyclopaedia Wikipedia provides free access to its content by receiving funding from charitable donations.[148]
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Notes
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Advertising is a marketing communication that employs an openly sponsored, non-personal message to promote or sell a product, service or idea.[1]:465 Sponsors of advertising are typically businesses wishing to promote their products or services. Advertising is differentiated from public relations in that an advertiser pays for and has control over the message. It differs from personal selling in that the message is non-personal, i.e., not directed to a particular individual.[1]:661,672
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Advertising is communicated through various mass media,[2] including traditional media such as newspapers, magazines, television, radio, outdoor advertising or direct mail; and new media such as search results, blogs, social media, websites or text messages. The actual presentation of the message in a medium is referred to as an advertisement, or "ad" or advert for short.
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Commercial ads often seek to generate increased consumption of their products or services through "branding", which associates a product name or image with certain qualities in the minds of consumers. On the other hand, ads that intend to elicit an immediate sale are known as direct-response advertising. Non-commercial entities that advertise more than consumer products or services include political parties, interest groups, religious organizations and governmental agencies. Non-profit organizations may use free modes of persuasion, such as a public service announcement. Advertising may also help to reassure employees or shareholders that a company is viable or successful.
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Modern advertising originated with the techniques introduced with tobacco advertising in the 1920s, most significantly with the campaigns of Edward Bernays, considered the founder of modern, "Madison Avenue" advertising.[3][4]
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Worldwide spending on advertising in 2015 amounted to an estimated US$529.43 billion.[5] Advertising's projected distribution for 2017 was 40.4% on TV, 33.3% on digital, 9% on newspapers, 6.9% on magazines, 5.8% on outdoor and 4.3% on radio.[6] Internationally, the largest ("Big Five") advertising-agency groups are Dentsu, Interpublic, Omnicom, Publicis, and WPP.[7]
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In Latin, advertere means "to turn towards".[8]
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Egyptians used papyrus to make sales messages and wall posters.[9] Commercial messages and political campaign displays have been found in the ruins of Pompeii and ancient Arabia. Lost and found advertising on papyrus was common in ancient Greece and ancient Rome. Wall or rock painting for commercial advertising is another manifestation of an ancient advertising form, which is present to this day in many parts of Asia, Africa, and South America. The tradition of wall painting can be traced back to Indian rock art paintings that date back to 4000 BC.[10]
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In ancient China, the earliest advertising known was oral, as recorded in the Classic of Poetry (11th to 7th centuries BC) of bamboo flutes played to sell confectionery. Advertisement usually takes in the form of calligraphic signboards and inked papers. A copper printing plate dated back to the Song dynasty used to print posters in the form of a square sheet of paper with a rabbit logo with "Jinan Liu's Fine Needle Shop" and "We buy high-quality steel rods and make fine-quality needles, to be ready for use at home in no time" written above and below[11] is considered the world's earliest identified printed advertising medium.[12]
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In Europe, as the towns and cities of the Middle Ages began to grow, and the general population was unable to read, instead of signs that read "cobbler", "miller", "tailor", or "blacksmith", images associated with their trade would be used such as a boot, a suit, a hat, a clock, a diamond, a horseshoe, a candle or even a bag of flour. Fruits and vegetables were sold in the city square from the backs of carts and wagons and their proprietors used street callers (town criers) to announce their whereabouts. The first compilation of such advertisements was gathered in "Les Crieries de Paris", a thirteenth-century poem by Guillaume de la Villeneuve.[13]
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In the 18th century advertisements started to appear in weekly newspapers in England. These early print advertisements were used mainly to promote books and newspapers, which became increasingly affordable with advances in the printing press; and medicines, which were increasingly sought after. However, false advertising and so-called "quack" advertisements became a problem, which ushered in the regulation of advertising content.
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Thomas J. Barratt of London has been called "the father of modern advertising".[14][15][16] Working for the Pears Soap company, Barratt created an effective advertising campaign for the company products, which involved the use of targeted slogans, images and phrases. One of his slogans, "Good morning. Have you used Pears' soap?" was famous in its day and into the 20th century.[17][18]
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Barratt introduced many of the crucial ideas that lie behind successful advertising and these were widely circulated in his day. He constantly stressed the importance of a strong and exclusive brand image for Pears and of emphasizing the product's availability through saturation campaigns. He also understood the importance of constantly reevaluating the market for changing tastes and mores, stating in 1907 that "tastes change, fashions change, and the advertiser has to change with them. An idea that was effective a generation ago would fall flat, stale, and unprofitable if presented to the public today. Not that the idea of today is always better than the older idea, but it is different – it hits the present taste."[15]
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As the economy expanded across the world during the 19th century, advertising grew alongside. In the United States, the success of this advertising format eventually led to the growth of mail-order advertising.
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In June 1836, French newspaper La Presse was the first to include paid advertising in its pages, allowing it to lower its price, extend its readership and increase its profitability and the formula was soon copied by all titles. Around 1840, Volney B. Palmer established the roots of the modern day advertising agency in Philadelphia. In 1842 Palmer bought large amounts of space in various newspapers at a discounted rate then resold the space at higher rates to advertisers. The actual ad – the copy, layout, and artwork – was still prepared by the company wishing to advertise; in effect, Palmer was a space broker. The situation changed when the first full-service advertising agency of N.W. Ayer & Son was founded in 1869 in Philadelphia. Ayer & Son offered to plan, create, and execute complete advertising campaigns for its customers. By 1900 the advertising agency had become the focal point of creative planning, and advertising was firmly established as a profession.
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[19] Around the same time, in France, Charles-Louis Havas extended the services of his news agency, Havas to include advertisement brokerage, making it the first French group to organize. At first, agencies were brokers for advertisement space in newspapers.[19]
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Advertising increased dramatically in the United States as industrialization expanded the supply of manufactured products. In 1919 it was 2.5 percent of gross domestic product (GDP) in the US, and it averaged 2.2 percent of GDP between then and at least 2007, though it may have declined dramatically since the Great Recession.
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Industry could not benefit from its increased productivity without a substantial increase in consumer spending. This contributed to the development of mass marketing designed to influence the population's economic behavior on a larger scale.[20] In the 1910s and 1920s, advertisers in the U.S. adopted the doctrine that human instincts could be targeted and harnessed – "sublimated" into the desire to purchase commodities.[21] Edward Bernays, a nephew of Sigmund Freud, became associated with the method and is sometimes called the founder of modern advertising and public relations.[22] Bernays claimed that:
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"[The] general principle, that men are very largely actuated by motives which they conceal from themselves, is as true of mass as of individual psychology. It is evident that the successful propagandist must understand the true motives and not be content to accept the reasons which men give for what they do."[23]
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In other words, selling products by appealing to the rational minds of customers (the main method used prior to Bernays) was much less effective than selling products based on the unconscious desires that Bernays felt were the true motivators of human action. "Sex sells" became a controversial issue, with techniques for titillating and enlarging the audience posing a challenge to conventional morality.[24][25]
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In the 1920s, under Secretary of Commerce Herbert Hoover, the American government promoted advertising. Hoover himself delivered an address to the Associated Advertising Clubs of the World in 1925 called 'Advertising Is a Vital Force in Our National Life."[26] In October 1929, the head of the U.S. Bureau of Foreign and Domestic Commerce, Julius Klein, stated "Advertising is the key to world prosperity."[27] This was part of the "unparalleled" collaboration between business and government in the 1920s, according to a 1933 European economic journal.[28]
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The tobacco companies became major advertisers in order to sell packaged cigarettes.[29] The tobacco companies pioneered the new advertising techniques when they hired Bernays to create positive associations with tobacco smoking.[3][4]
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Advertising was also used as a vehicle for cultural assimilation, encouraging workers to exchange their traditional habits and community structure in favor of a shared "modern" lifestyle.[30] An important tool for influencing immigrant workers was the American Association of Foreign Language Newspapers (AAFLN). The AAFLN was primarily an advertising agency but also gained heavily centralized control over much of the immigrant press.[31][32]
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At the turn of the 20th century, advertising was one of the few career choices for women. Since women were responsible for most household purchasing done, advertisers and agencies recognized the value of women's insight during the creative process. In fact, the first American advertising to use a sexual sell was created by a woman – for a soap product. Although tame by today's standards, the advertisement featured a couple with the message "A skin you love to touch".[33]
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In the 1920s psychologists Walter D. Scott and John B. Watson contributed applied psychological theory to the field of advertising. Scott said, "Man has been called the reasoning animal but he could with greater truthfulness be called the creature of suggestion. He is reasonable, but he is to a greater extent suggestible".[34] He demonstrated this through his advertising technique of a direct command to the consumer.
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In the early 1920s, the first radio stations were established by radio equipment manufacturers, followed by non-profit organizations such as schools, clubs and civic groups who also set up their own stations.[35] Retailer and consumer goods manufacturers quickly recognised radio's potential to reach consumers in their home and soon adopted advertising techniques that would allow their messages to stand out; slogans, mascots, and jingles began to appear on radio in the 1920s and early television in the 1930s.[36]
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The rise of mass media communications allowed manufacturers of branded goods to bypass retailers by advertising directly to consumers. This was a major paradigm shift which forced manufacturers to focus on the brand and stimulated the need for superior insights into consumer purchasing, consumption and usage behaviour; their needs, wants and aspirations.[37] The earliest radio drama series were sponsored by soap manufacturers and the genre became known as a soap opera.[38] Before long, radio station owners realized they could increase advertising revenue by selling 'air-time' in small time allocations which could be sold to multiple businesses. By the 1930s, these advertising spots, as the packets of time became known, were being sold by the station's geographical sales representatives, ushering in an era of national radio advertising.[39]
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By the 1940s, manufacturers began to recognize the way in which consumers were developing personal relationships with their brands in a social/psychological/anthropological sense.[40] Advertisers began to use motivational research and consumer research to gather insights into consumer purchasing. Strong branded campaigns for Chrysler and Exxon/Esso, using insights drawn research methods from psychology and cultural anthropology, led to some of the most enduring campaigns of the 20th-century.[41]
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In the early 1950s, the DuMont Television Network began the modern practice of selling advertisement time to multiple sponsors. Previously, DuMont had trouble finding sponsors for many of their programs and compensated by selling smaller blocks of advertising time to several businesses. This eventually became the standard for the commercial television industry in the United States. However, it was still a common practice to have single sponsor shows, such as The United States Steel Hour. In some instances the sponsors exercised great control over the content of the show – up to and including having one's advertising agency actually writing the show.[citation needed] The single sponsor model is much less prevalent now, a notable exception being the Hallmark Hall of Fame.[citation needed]
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The late 1980s and early 1990s saw the introduction of cable television and particularly MTV. Pioneering the concept of the music video, MTV ushered in a new type of advertising: the consumer tunes in for the advertising message, rather than it being a by-product or afterthought. As cable and satellite television became increasingly prevalent, specialty channels emerged, including channels entirely devoted to advertising, such as QVC, Home Shopping Network, and ShopTV Canada.[42]
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With the advent of the ad server, online advertising grew, contributing to the "dot-com" boom of the 1990s.[43] Entire corporations operated solely on advertising revenue, offering everything from coupons to free Internet access. At the turn of the 21st century, some websites, including the search engine Google, changed online advertising by personalizing ads based on web browsing behavior. This has led to other similar efforts and an increase in interactive advertising.[44]
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The share of advertising spending relative to GDP has changed little across large changes in media since 1925. In 1925, the main advertising media in America were newspapers, magazines, signs on streetcars, and outdoor posters. Advertising spending as a share of GDP was about 2.9 percent. By 1998, television and radio had become major advertising media; by 2017, the balance between broadcast and online advertising had shifted, with online spending exceeding broadcast.[45] Nonetheless, advertising spending as a share of GDP was slightly lower – about 2.4 percent.[46]
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Guerrilla marketing involves unusual approaches such as staged encounters in public places, giveaways of products such as cars that are covered with brand messages, and interactive advertising where the viewer can respond to become part of the advertising message. This type of advertising is unpredictable, which causes consumers to buy the product or idea.[47] This reflects an increasing trend of interactive and "embedded" ads, such as via product placement, having consumers vote through text messages, and various campaigns utilizing social network services such as Facebook or Twitter.[48]
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The advertising business model has also been adapted in recent years.[when?][clarification needed] In media for equity, advertising is not sold, but provided to start-up companies in return for equity. If the company grows and is sold, the media companies receive cash for their shares.
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Domain name registrants (usually those who register and renew domains as an investment) sometimes "park" their domains and allow advertising companies to place ads on their sites in return for per-click payments. These ads are typically driven by pay per click search engines like Google or Yahoo, but ads can sometimes be placed directly on targeted domain names through a domain lease or by making contact with the registrant of a domain name that describes a product. Domain name registrants are generally easy to identify through WHOIS records that are publicly available at registrar websites.[49]
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Advertising may be categorized in a variety of ways, including by style, target audience, geographic scope, medium, or purpose.[2]:9–15 For example, in print advertising, classification by style can include display advertising (ads with design elements sold by size) vs. classified advertising (ads without design elements sold by the word or line). Advertising may be local, national or global. An ad campaign may be directed toward consumers or to businesses. The purpose of an ad may be to raise awareness (brand advertising), or to elicit an immediate sale (direct response advertising). The term above the line (ATL) is used for advertising involving mass media; more targeted forms of advertising and promotion are referred to as below the line (BTL).[50][51] The two terms date back to 1954 when Procter & Gamble began paying their advertising agencies differently from other promotional agencies.[52] In the 2010s, as advertising technology developed, a new term, through the line (TTL) began to come into use, referring to integrated advertising campaigns.[53][54]
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Virtually any medium can be used for advertising. Commercial advertising media can include wall paintings, billboards, street furniture components, printed flyers and rack cards, radio, cinema and television adverts, web banners, mobile telephone screens, shopping carts, web popups, skywriting, bus stop benches, human billboards and forehead advertising, magazines, newspapers, town criers, sides of buses, banners attached to or sides of airplanes ("logojets"), in-flight advertisements on seatback tray tables or overhead storage bins, taxicab doors, roof mounts and passenger screens, musical stage shows, subway platforms and trains, elastic bands on disposable diapers, doors of bathroom stalls, stickers on apples in supermarkets, shopping cart handles (grabertising), the opening section of streaming audio and video, posters, and the backs of event tickets and supermarket receipts. Any situation in which an "identified" sponsor pays to deliver their message through a medium is advertising.[55]
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A new advertising approach is known as advanced advertising, which is data-driven advertising, using large quantities of data, precise measuring tools and precise targeting.[66] Advanced advertising also makes it easier for companies which sell ad-space to attribute customer purchases to the ads they display or broadcast.[67]
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Increasingly, other media are overtaking many of the "traditional" media such as television, radio and newspaper because of a shift toward the usage of the Internet for news and music as well as devices like digital video recorders (DVRs) such as TiVo.[68]
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Online advertising began with unsolicited bulk e-mail advertising known as "e-mail spam". Spam has been a problem for e-mail users since 1978.[69] As new online communication channels became available, advertising followed. The first banner ad appeared on the World Wide Web in 1994.[70] Prices of Web-based advertising space are dependent on the "relevance" of the surrounding web content and the traffic that the website receives.[citation needed]
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In online display advertising, display ads generate awareness quickly. Unlike search, which requires someone to be aware of a need, display advertising can drive awareness of something new and without previous knowledge. Display works well for direct response. Display is not only used for generating awareness, it's used for direct response campaigns that link to a landing page with a clear 'call to action'.[71]
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As the mobile phone became a new mass medium in 1998 when the first paid downloadable content appeared on mobile phones in Finland,[72][citation needed] mobile advertising followed, also first launched in Finland in 2000.[citation needed] By 2007 the value of mobile advertising had reached $2 billion and providers such as Admob delivered billions of mobile ads.[citation needed]
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More advanced mobile ads include banner ads, coupons, Multimedia Messaging Service picture and video messages, advergames and various engagement marketing campaigns. A particular feature driving mobile ads is the 2D barcode, which replaces the need to do any typing of web addresses, and uses the camera feature of modern phones to gain immediate access to web content. 83 percent of Japanese mobile phone users already are active users of 2D barcodes.[citation needed]
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Some companies have proposed placing messages or corporate logos on the side of booster rockets and the International Space Station.[citation needed]
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Unpaid advertising (also called "publicity advertising"), can include personal recommendations ("bring a friend", "sell it"), spreading buzz, or achieving the feat of equating a brand with a common noun (in the United States, "Xerox" = "photocopier", "Kleenex" = tissue, "Vaseline" = petroleum jelly, "Hoover" = vacuum cleaner, and "Band-Aid" = adhesive bandage). However, some companies[which?] oppose the use of their brand name to label an object. Equating a brand with a common noun also risks turning that brand into a generic trademark – turning it into a generic term which means that its legal protection as a trademark is lost.[73][disputed – discuss]
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From time to time, The CW Television Network airs short programming breaks called "Content Wraps", to advertise one company's product during an entire commercial break. The CW pioneered "content wraps" and some products featured were Herbal Essences, Crest, Guitar Hero II, CoverGirl, and Toyota.[74][75]
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A new promotion concept has appeared, "ARvertising", advertising on augmented reality technology.[76]
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Controversy exists on the effectiveness of subliminal advertising (see mind control), and the pervasiveness of mass messages (propaganda).
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With the Internet came many new advertising opportunities. Pop-up, Flash, banner, pop-under, advergaming, and email advertisements (all of which are often unwanted or spam in the case of email) are now commonplace. Particularly since the rise of "entertaining" advertising, some people may like an advertisement enough to wish to watch it later or show a friend.[citation needed] In general, the advertising community has not yet made this easy, although some have used the Internet to widely distribute their ads to anyone willing to see or hear them. In the last three quarters of 2009, mobile and Internet advertising grew by 18% and 9% respectively, while older media advertising saw declines: −10.1% (TV), −11.7% (radio), −14.8% (magazines) and −18.7% (newspapers).[citation needed] Between 2008 and 2014, U.S. newspapers lost more than half their print advertising revenue.[78]
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Another significant trend regarding future of advertising is the growing importance of the niche market using niche or targeted ads. Also brought about by the Internet and the theory of the long tail, advertisers will have an increasing ability to reach specific audiences. In the past, the most efficient way to deliver a message was to blanket the largest mass market audience possible.[citation needed] However, usage tracking, customer profiles and the growing popularity of niche content brought about by everything from blogs to social networking sites, provide advertisers with audiences that are smaller but much better defined,[citation needed] leading to ads that are more relevant to viewers and more effective for companies' marketing products. Among others, Comcast Spotlight is one such advertiser employing this method in their video on demand menus. These advertisements are targeted to a specific group and can be viewed by anyone wishing to find out more about a particular business or practice, from their home. This causes the viewer to become proactive and actually choose what advertisements they want to view.[79]
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Niche marketing could also be helped by bringing the issue of colour into advertisements. Different colours play major roles when it comes to marketing strategies, for example, seeing the blue can promote a sense of calmness and gives a sense of security which is why many social networks such as Facebook use blue in their logos.
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Google AdSense is an example of niche marketing. Google calculates the primary purpose of a website and adjusts ads accordingly; it uses keywords on the page (or even in emails) to find the general ideas of topics disused and places ads that will most likely be clicked on by viewers of the email account or website visitors.
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The concept of crowdsourcing has given way to the trend of user-generated advertisements. User-generated ads are created by people, as opposed to an advertising agency or the company themselves, often resulting from brand sponsored advertising competitions. For the 2007 Super Bowl, the Frito-Lays division of PepsiCo held the "Crash the Super Bowl" contest, allowing people to create their own Doritos commercials.[80] Chevrolet held a similar competition for their Tahoe line of SUVs.[80] Due to the success of the Doritos user-generated ads in the 2007 Super Bowl, Frito-Lays relaunched the competition for the 2009 and 2010 Super Bowl. The resulting ads were among the most-watched and most-liked Super Bowl ads. In fact, the winning ad that aired in the 2009 Super Bowl was ranked by the USA Today Super Bowl Ad Meter as the top ad for the year while the winning ads that aired in the 2010 Super Bowl were found by Nielsen's BuzzMetrics to be the "most buzzed-about".[81][82] Another example of companies using crowdsourcing successfully is the beverage company Jones Soda that encourages consumers to participate in the label design themselves.[83]
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This trend has given rise to several online platforms that host user-generated advertising competitions on behalf of a company. Founded in 2007, Zooppa has launched ad competitions for brands such as Google, Nike, Hershey's, General Mills, Microsoft, NBC Universal, Zinio, and Mini Cooper.[84] Crowdsourcing remains controversial, as the long-term impact on the advertising industry is still unclear.[85]
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Advertising has gone through five major stages of development: domestic, export, international, multi-national, and global. For global advertisers, there are four, potentially competing, business objectives that must be balanced when developing worldwide advertising: building a brand while speaking with one voice, developing economies of scale in the creative process, maximising local effectiveness of ads, and increasing the company's speed of implementation. Born from the evolutionary stages of global marketing are the three primary and fundamentally different approaches to the development of global advertising executions: exporting executions, producing local executions, and importing ideas that travel.[86]
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Advertising research is key to determining the success of an ad in any country or region. The ability to identify which elements and/or moments of an ad contribute to its success is how economies of scale are maximized. Once one knows what works in an ad, that idea or ideas can be imported by any other market. Market research measures, such as Flow of Attention, Flow of Emotion and branding moments provide insight into what is working in an ad in any country or region because the measures are based on the visual, not verbal, elements of the ad.[87]
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Foreign governments,[which?] particularly those that own marketable commercial products or services, often promote their interests and positions through the advertising of those goods because the target audience is not only largely unaware of the forum as a vehicle for foreign messaging but also willing to receive the message while in a mental state of absorbing information from advertisements during television commercial breaks, while reading a periodical, or while passing by billboards in public spaces. A prime example of this messaging technique is advertising campaigns to promote international travel. While advertising foreign destinations and services may stem from the typical goal of increasing revenue by drawing more tourism, some travel campaigns carry the additional or alternative intended purpose of promoting good sentiments or improving existing ones among the target audience towards a given nation or region. It is common for advertising promoting foreign countries to be produced and distributed by the tourism ministries of those countries, so these ads often carry political statements and/or depictions of the foreign government's desired international public perception. Additionally, a wide range of foreign airlines and travel-related services which advertise separately from the destinations, themselves, are owned by their respective governments; examples include, though are not limited to, the Emirates airline (Dubai), Singapore Airlines (Singapore), Qatar Airways (Qatar), China Airlines (Taiwan/Republic of China), and Air China (People's Republic of China). By depicting their destinations, airlines, and other services in a favorable and pleasant light, countries market themselves to populations abroad in a manner that could mitigate prior public impressions.
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In the realm of advertising agencies, continued industry diversification has seen observers note that "big global clients don't need big global agencies any more".[88] This is reflected by the growth of non-traditional agencies in various global markets, such as Canadian business TAXI and SMART in Australia and has been referred to as "a revolution in the ad world".[89]
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The ability to record shows on digital video recorders (such as TiVo) allow watchers to record the programs for later viewing, enabling them to fast forward through commercials. Additionally, as more seasons of pre-recorded box sets are offered for sale of television programs; fewer people watch the shows on TV. However, the fact that these sets are sold, means the company will receive additional profits from these sets.
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To counter this effect, a variety of strategies have been employed. Many advertisers have opted for product placement on TV shows like Survivor. Other strategies include integrating advertising with internet-connected program guidess (EPGs), advertising on companion devices (like smartphones and tablets) during the show, and creating mobile apps for TV programs. Additionally, some like brands have opted for social television sponsorship.[90]
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The emerging technology of drone displays has recently been used for advertising purposes.[91]
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In recent years there have been several media literacy initiatives, and more specifically concerning advertising, that seek to empower citizens in the face of media advertising campaigns.[92]
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Advertising education has become popular with bachelor, master and doctorate degrees becoming available in the emphasis.[citation needed] A surge in advertising interest is typically attributed to the strong relationship advertising plays in cultural and technological changes, such as the advance of online social networking.[citation needed] A unique model for teaching advertising is the student-run advertising agency, where advertising students create campaigns for real companies.[93] Organizations such as the American Advertising Federation establish companies with students to create these campaigns.[citation needed]
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Advertising is at the front of delivering the proper message to customers and prospective customers. The purpose of advertising is to inform the consumers about their product and convince customers that a company's services or products are the best, enhance the image of the company, point out and create a need for products or services, demonstrate new uses for established products, announce new products and programs, reinforce the salespeople's individual messages, draw customers to the business, and to hold existing customers.[94]
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Sales promotions are another way to advertise. Sales promotions are double purposed because they are used to gather information about what type of customers one draws in and where they are, and to jump start sales. Sales promotions include things like contests and games, sweepstakes, product giveaways, samples coupons, loyalty programs, and discounts. The ultimate goal of sales promotions is to stimulate potential customers to action.[95]
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While advertising can be seen as necessary for economic growth,[27] it is not without social costs. Unsolicited commercial e-mail and other forms of spam have become so prevalent as to have become a major nuisance to users of these services, as well as being a financial burden on internet service providers.[96] Advertising is increasingly invading public spaces, such as schools, which some critics argue is a form of child exploitation.[97] This increasing difficulty in limiting exposure to specific audiences can result in negative backlash for advertisers.[98] In tandem with these criticisms, the advertising industry has seen low approval rates in surveys and negative cultural portrayals.[99]
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One of the most controversial criticisms of advertisement in the present day is that of the predominance of advertising of foods high in sugar, fat, and salt specifically to children. Critics claim that food advertisements targeting children are exploitive and are not sufficiently balanced with proper nutritional education to help children understand the consequences of their food choices. Additionally, children may not understand that they are being sold something, and are therefore more impressionable.[100] Michelle Obama has criticized large food companies for advertising unhealthy foods largely towards children and has requested that food companies either limit their advertising to children or advertise foods that are more in line with dietary guidelines.[101] The other criticisms include the change that are brought by those advertisements on the society and also the deceiving ads that are aired and published by the corporations. Cosmetic and health industry are the ones which exploited the highest and created reasons of concern.[102]
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There have been increasing efforts to protect the public interest by regulating the content and the influence of advertising. Some examples include restrictions for advertising alcohol, tobacco or gambling imposed in many countries, as well as the bans around advertising to children, which exist in parts of Europe. Advertising regulation focuses heavily on the veracity of the claims and as such, there are often tighter restrictions placed around advertisements for food and healthcare products.[103]
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The advertising industries within some countries rely less on laws and more on systems of self-regulation.[103][104][105] Advertisers and the media agree on a code of advertising standards that they attempt to uphold. The general aim of such codes is to ensure that any advertising is 'legal, decent, honest and truthful'. Some self-regulatory organizations are funded by the industry, but remain independent, with the intent of upholding the standards or codes like the Advertising Standards Authority in the UK.[106]
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In the UK, most forms of outdoor advertising such as the display of billboards is regulated by the UK Town and County Planning system. Currently, the display of an advertisement without consent from the Planning Authority is a criminal offense liable to a fine of £2,500 per offense.[107] In the US, many communities believe that many forms of outdoor advertising blight the public realm.[108] As long ago as the 1960s in the US, there were attempts to ban billboard advertising in the open countryside.[109] Cities such as São Paulo have introduced an outright ban[110] with London also having specific legislation to control unlawful displays.
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Some governments restrict the languages that can be used in advertisements, but advertisers may employ tricks to try avoiding them. In France for instance, advertisers sometimes print English words in bold and French translations in fine print to deal with Article 120 of the 1994 Toubon Law limiting the use of English.[111]
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The advertising of pricing information is another topic of concern for governments. In the United States for instance, it is common for businesses to only mention the existence and amount of applicable taxes at a later stage of a transaction.[112] In Canada and New Zealand, taxes can be listed as separate items, as long as they are quoted up-front.[113][114] In most other countries, the advertised price must include all applicable taxes, enabling customers to easily know how much it will cost them.[115][116][117]
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Various competing models of hierarchies of effects attempt to provide a theoretical underpinning to advertising practice.[clarification needed][118]
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The marketing mix was proposed by professor E. Jerome McCarthy in the 1960s.[122] It consists of four basic elements called the "four Ps". Product is the first P representing the actual product. Price represents the process of determining the value of a product. Place represents the variables of getting the product to the consumer such as distribution channels, market coverage and movement organization. The last P stands for Promotion which is the process of reaching the target market and convincing them to buy the product.
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In the 1990s, the concept of four Cs was introduced as a more customer-driven replacement of four P's.[123] There are two theories based on four Cs: Lauterborn's four Cs (consumer, cost, communication, convenience)
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[124] and Shimizu's four Cs (commodity, cost, communication, channel) in the 7Cs Compass Model (Co-marketing). Communications can include advertising, sales promotion, public relations, publicity, personal selling, corporate identity, internal communication, SNS, and MIS.[125][126][127][128]
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Advertising research is a specialized form of research that works to improve the effectiveness and efficiency of advertising. It entails numerous forms of research which employ different methodologies. Advertising research includes pre-testing (also known as copy testing) and post-testing of ads and/or campaigns.
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Pre-testing includes a wide range of qualitative and quantitative techniques, including: focus groups, in-depth target audience interviews (one-on-one interviews), small-scale quantitative studies and physiological measurement. The goal of these investigations is to better understand how different groups respond to various messages and visual prompts, thereby providing an assessment of how well the advertisement meets its communications goals.[129]
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Post-testing employs many of the same techniques as pre-testing, usually with a focus on understanding the change in awareness or attitude attributable to the advertisement. With the emergence of digital advertising technologies, many firms have begun to continuously post-test ads using real-time data. This may take the form of A/B split-testing or multivariate testing.
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Continuous ad tracking and the Communicus System are competing examples of post-testing advertising research types.[130]
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Meanings between consumers and marketers depict signs and symbols that are encoded in everyday objects.[131] Semiotics is the study of signs and how they are interpreted. Advertising has many hidden signs and meanings within brand names, logos, package designs, print advertisements, and television advertisements. Semiotics aims to study and interpret the message being conveyed in (for example) advertisements. Logos and advertisements can be interpreted at two levels – known as the surface level and the underlying level. The surface level uses signs creatively to create an image or personality for a product.[citation needed] These signs can be images, words, fonts, colors, or slogans. The underlying level is made up of hidden meanings. The combination of images, words, colors, and slogans must be interpreted by the audience or consumer.[132] The "key to advertising analysis" is the signifier and the signified. The signifier is the object and the signified is the mental concept.[133] A product has a signifier and a signified. The signifier is the color, brand name, logo design, and technology. The signified has two meanings known as denotative and connotative. The denotative meaning is the meaning of the product. A television's denotative meaning might be that it is high definition. The connotative meaning is the product's deep and hidden meaning. A connotative meaning of a television would be that it is top-of-the-line.[134]
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Apple's commercials[when?] used a black silhouette of a person that was the age of Apple's target market. They placed the silhouette in front of a blue screen so that the picture behind the silhouette could be constantly changing. However, the one thing that stays the same in these ads is that there is music in the background and the silhouette is listening to that music on a white iPod through white headphones. Through advertising, the white color on a set of earphones now signifies that the music device is an iPod. The white color signifies almost all of Apple's products.[135]
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The semiotics of gender plays a key influence on the way in which signs are interpreted. When considering gender roles in advertising, individuals are influenced by three categories. Certain characteristics of stimuli may enhance or decrease the elaboration of the message (if the product is perceived as feminine or masculine). Second, the characteristics of individuals can affect attention and elaboration of the message (traditional or non-traditional gender role orientation). Lastly, situational factors may be important to influence the elaboration of the message.[136]
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There are two types of marketing communication claims-objective and subjective.[137] Objective claims stem from the extent to which the claim associates the brand with a tangible product or service feature. For instance, a camera may have auto-focus features. Subjective claims convey emotional, subjective, impressions of intangible aspects of a product or service. They are non-physical features of a product or service that cannot be directly perceived, as they have no physical reality. For instance the brochure has a beautiful design.[138] Males tend to respond better to objective marketing-communications claims while females tend to respond better to subjective marketing communications claims.[139]
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Voiceovers are commonly used in advertising. Most voiceovers are done by men, with figures of up to 94% having been reported.[140] There have been more female voiceovers in recent years,[when?] but mainly for food, household products, and feminine-care products.[141]
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+
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According to a 1977 study by David Statt, females process information comprehensively, while males process information through heuristic devices such as procedures, methods or strategies for solving problems, which could have an effect on how they interpret advertising.[142][need quotation to verify] According to this study, men prefer to have available and apparent cues to interpret the message, whereas females engage in more creative, associative, imagery-laced interpretation. Later research by a Danish team[143] found that advertising attempts to persuade men to improve their appearance or performance, whereas its approach to women aims at transformation toward an impossible ideal of female presentation. In Paul Suggett's article "The Objectification of Women in Advertising"[144] he discusses the negative impact that these women in advertisements, who are too perfect to be real, have on women, as well as men, in real life. Advertising's manipulation of women's aspiration to these ideal types as portrayed in film, in erotic art, in advertising, on stage, within music videos, and through other media exposures requires at least a conditioned rejection of female reality, and thereby takes on a highly ideological cast. Studies show that these expectations of women and young girls negatively impact their views about their bodies and appearances. These advertisements are directed towards men. Not everyone agrees: one critic viewed this monologic, gender-specific interpretation of advertising as excessively skewed and politicized.[145][need quotation to verify] There are some companies, however, like Dove and Aerie that are creating commercials to portray more natural women, with less post production manipulation, so more women and young girls are able to relate to them.[citation needed]
|
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+
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More recent research by Martin (2003) reveals that males and females differ in how they react to advertising depending on their mood at the time of exposure to the ads, and on the affective tone of the advertising. When feeling sad, males prefer happy ads to boost their mood. In contrast, females prefer happy ads when they are feeling happy. The television programs in which ads are embedded influence a viewer's mood state.[146] Susan Wojcicki, author of the article "Ads that Empower Women don’t just Break Stereotypes—They’re also Effective"[147] discusses how advertising to women has changed since the first Barbie commercial where a little girl tells the doll that, she wants to be just like her. Little girls grow up watching advertisements of scantily clad women advertising things from trucks to burgers, and Wojcicki states that this shows girls that they are either arm candy or eye candy.
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+
Other approaches to revenue include donations, paid subscriptions and microtransactions. Websites and applications are "ad-free" when not using ads at all for revenue. For example, the online encyclopaedia Wikipedia provides free access to its content by receiving funding from charitable donations.[148]
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Notes
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1 |
+
A pulsar (from pulse and -ar as in quasar)[1] is a highly magnetized rotating neutron star that emits beams of electromagnetic radiation out of its magnetic poles.[2] This radiation can be observed only when a beam of emission is pointing toward Earth (much like the way a lighthouse can be seen only when the light is pointed in the direction of an observer), and is responsible for the pulsed appearance of emission. Neutron stars are very dense, and have short, regular rotational periods. This produces a very precise interval between pulses that ranges from milliseconds to seconds for an individual pulsar. Pulsars are one of the candidates for the source of ultra-high-energy cosmic rays (see also centrifugal mechanism of acceleration).
|
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+
|
3 |
+
The periods of pulsars make them very useful tools for astronomers. Observations of a pulsar in a binary neutron star system were used to indirectly confirm the existence of gravitational radiation. The first extrasolar planets were discovered around a pulsar, PSR B1257+12. In 1983, certain types of pulsars exceeded atomic clocks in their accuracy in keeping time.[3]
|
4 |
+
|
5 |
+
The first pulsar was observed on November 28, 1967, by Jocelyn Bell Burnell and Antony Hewish.[4][5][6] They observed pulses separated by 1.33 seconds that originated from the same location in the sky, and kept to sidereal time. In looking for explanations for the pulses, the short period of the pulses eliminated most astrophysical sources of radiation, such as stars, and since the pulses followed sidereal time, it could not be human-made radio frequency interference.
|
6 |
+
|
7 |
+
When observations with another telescope confirmed the emission, it eliminated any sort of instrumental effects. At this point, Bell Burnell said of herself and Hewish that "we did not really believe that we had picked up signals from another civilization, but obviously the idea had crossed our minds and we had no proof that it was an entirely natural radio emission. It is an interesting problem—if one thinks one may have detected life elsewhere in the universe, how does one announce the results responsibly?"[7] Even so, they nicknamed the signal LGM-1, for "little green men" (a playful name for intelligent beings of extraterrestrial origin).
|
8 |
+
|
9 |
+
It was not until a second pulsating source was discovered in a different part of the sky that the "LGM hypothesis" was entirely abandoned.[8] Their pulsar was later dubbed CP 1919, and is now known by a number of designators including PSR 1919+21 and PSR J1921+2153. Although CP 1919 emits in radio wavelengths, pulsars have subsequently been found to emit in visible light, X-ray, and gamma ray wavelengths.[9]
|
10 |
+
The word "pulsar" is a portmanteau of 'pulsating' and 'quasar', and first appeared in print in 1968:
|
11 |
+
|
12 |
+
An entirely novel kind of star came to light on Aug. 6 last year and was referred to, by astronomers, as LGM (Little Green Men). Now it is thought to be a novel type between a white dwarf and a neutron [star]. The name Pulsar is likely to be given to it. Dr. A. Hewish told me yesterday: '... I am sure that today every radio telescope is looking at the Pulsars.'[10]
|
13 |
+
|
14 |
+
The existence of neutron stars was first proposed by Walter Baade and Fritz Zwicky in 1934, when they argued that a small, dense star consisting primarily of neutrons would result from a supernova.[11] Based on the idea of magnetic flux conservation from magnetic main sequence stars, Lodewijk Woltjer proposed in 1964 that such neutron stars might contain magnetic fields as large as 1014 to 1016 G.[12] In 1967, shortly before the discovery of pulsars, Franco Pacini suggested that a rotating neutron star with a magnetic field would emit radiation, and even noted that such energy could be pumped into a supernova remnant around a neutron star, such as the Crab Nebula.[13] After the discovery of the first pulsar, Thomas Gold independently suggested a rotating neutron star model similar to that of Pacini, and explicitly argued that this model could explain the pulsed radiation observed by Bell Burnell and Hewish.[14] The discovery of the Crab pulsar later in 1968 seemed to provide confirmation of the rotating neutron star model of pulsars. The Crab pulsar has a 33-millisecond pulse period, which was too short to be consistent with other proposed models for pulsar emission. Moreover, the Crab pulsar is so named because it is located at the center of the Crab Nebula, consistent with the 1933 prediction of Baade and Zwicky.[15]
|
15 |
+
|
16 |
+
In 1974, Antony Hewish and Martin Ryle, who had developed revolutionary radio telescopes, became the first astronomers to be awarded the Nobel Prize in Physics, with the Royal Swedish Academy of Sciences noting that Hewish played a "decisive role in the discovery of pulsars".[16] Considerable controversy is associated with the fact that Hewish was awarded the prize while Bell, who made the initial discovery while she was his PhD student, was not. Bell claims no bitterness upon this point, supporting the decision of the Nobel prize committee.[17]
|
17 |
+
|
18 |
+
In 1974, Joseph Hooton Taylor, Jr. and Russell Hulse discovered for the first time a pulsar in a binary system, PSR B1913+16. This pulsar orbits another neutron star with an orbital period of just eight hours. Einstein's theory of general relativity predicts that this system should emit strong gravitational radiation, causing the orbit to continually contract as it loses orbital energy. Observations of the pulsar soon confirmed this prediction, providing the first ever evidence of the existence of gravitational waves. As of 2010, observations of this pulsar continue to agree with general relativity.[18] In 1993, the Nobel Prize in Physics was awarded to Taylor and Hulse for the discovery of this pulsar.[19]
|
19 |
+
|
20 |
+
In 1982, Don Backer led a group which discovered PSR B1937+21, a pulsar with a rotation period of just 1.6 milliseconds (38,500 rpm).[20] Observations soon revealed that its magnetic field was much weaker than ordinary pulsars, while further discoveries cemented the idea that a new class of object, the "millisecond pulsars" (MSPs) had been found. MSPs are believed to be the end product of X-ray binaries. Owing to their extraordinarily rapid and stable rotation, MSPs can be used by astronomers as clocks rivaling the stability of the best atomic clocks on Earth. Factors affecting the arrival time of pulses at Earth by more than a few hundred nanoseconds can be easily detected and used to make precise measurements. Physical parameters accessible through pulsar timing include the 3D position of the pulsar, its proper motion, the electron content of the interstellar medium along the propagation path, the orbital parameters of any binary companion, the pulsar rotation period and its evolution with time. (These are computed from the raw timing data by Tempo, a computer program specialized for this task.) After these factors have been taken into account, deviations between the observed arrival times and predictions made using these parameters can be found and attributed to one of three possibilities: intrinsic variations in the spin period of the pulsar, errors in the realization of Terrestrial Time against which arrival times were measured, or the presence of background gravitational waves. Scientists are currently attempting to resolve these possibilities by comparing the deviations seen between several different pulsars, forming what is known as a pulsar timing array. The goal of these efforts is to develop a pulsar-based time standard precise enough to make the first ever direct detection of gravitational waves.
|
21 |
+
In June 2006, the astronomer John Middleditch and his team at LANL announced the first prediction of pulsar glitches with observational data from the Rossi X-ray Timing Explorer. They used observations of the pulsar PSR J0537-6910.
|
22 |
+
|
23 |
+
In 1992, Aleksander Wolszczan discovered the first extrasolar planets around PSR B1257+12. This discovery presented important evidence concerning the widespread existence of planets outside the Solar System, although it is very unlikely that any life form could survive in the environment of intense radiation near a pulsar.
|
24 |
+
|
25 |
+
In 2016, AR Scorpii was identified as the first pulsar in which the compact object is a white dwarf instead of a neutron star.[21] Because its moment of inertia is much higher than that of a neutron star, the white dwarf in this system rotates once every 1.97 minutes, far slower than neutron-star pulsars.[22] The system displays strong pulsations from ultraviolet to radio wavelengths, powered by the spin-down of the strongly magnetized white dwarf.[21]
|
26 |
+
|
27 |
+
Initially pulsars were named with letters of the discovering observatory followed by their right ascension (e.g. CP 1919). As more pulsars were discovered, the letter code became unwieldy, and so the convention then arose of using the letters PSR (Pulsating Source of Radio) followed by the pulsar's right ascension and degrees of declination (e.g. PSR 0531+21) and sometimes declination to a tenth of a degree (e.g. PSR 1913+16.7). Pulsars appearing very close together sometimes have letters appended (e.g. PSR 0021−72C and PSR 0021−72D).
|
28 |
+
|
29 |
+
The modern convention prefixes the older numbers with a B (e.g. PSR B1919+21), with the B meaning the coordinates are for the 1950.0 epoch. All new pulsars have a J indicating 2000.0 coordinates and also have declination including minutes (e.g. PSR J1921+2153). Pulsars that were discovered before 1993 tend to retain their B names rather than use their J names (e.g. PSR J1921+2153 is more commonly known as PSR B1919+21). Recently discovered pulsars only have a J name (e.g. PSR J0437−4715). All pulsars have a J name that provides more precise coordinates of its location in the sky.[23]
|
30 |
+
|
31 |
+
The events leading to the formation of a pulsar begin when the core of a massive star is compressed during a supernova, which collapses into a neutron star. The neutron star retains most of its angular momentum, and since it has only a tiny fraction of its progenitor's radius (and therefore its moment of inertia is sharply reduced), it is formed with very high rotation speed. A beam of radiation is emitted along the magnetic axis of the pulsar, which spins along with the rotation of the neutron star. The magnetic axis of the pulsar determines the direction of the electromagnetic beam, with the magnetic axis not necessarily being the same as its rotational axis. This misalignment causes the beam to be seen once for every rotation of the neutron star, which leads to the "pulsed" nature of its appearance.
|
32 |
+
|
33 |
+
In rotation-powered pulsars, the beam originates from the rotational energy of the neutron star, which generates an electrical field from the movement of the very strong magnetic field, resulting in the acceleration of protons and electrons on the star surface and the creation of an electromagnetic beam emanating from the poles of the magnetic field.[24][25] Observations by NICER of J0030-0451, both beams originate from hotspots located on the south pole and that there may be more than two such hotspots on that star.[26][27] This rotation slows down over time as electromagnetic power is emitted. When a pulsar's spin period slows down sufficiently, the radio pulsar mechanism is believed to turn off (the so-called "death line"). This turn-off seems to take place after about 10–100 million years, which means of all the neutron stars born in the 13.6 billion year age of the universe, around 99% no longer pulsate.[28]
|
34 |
+
|
35 |
+
Though the general picture of pulsars as rapidly rotating neutron stars is widely accepted, Werner Becker of the Max Planck Institute for Extraterrestrial Physics said in 2006, "The theory of how pulsars emit their radiation is still in its infancy, even after nearly forty years of work."[29]
|
36 |
+
|
37 |
+
Three distinct classes of pulsars are currently known to astronomers, according to the source of the power of the electromagnetic radiation:
|
38 |
+
|
39 |
+
Although all three classes of objects are neutron stars, their observable behavior and the underlying physics are quite different. There are, however, connections. For example, X-ray pulsars are probably old rotationally-powered pulsars that have already lost most of their power, and have only become visible again after their binary companions had expanded and began transferring matter on to the neutron star. The process of accretion can in turn transfer enough angular momentum to the neutron star to "recycle" it as a rotation-powered millisecond pulsar. As this matter lands on the neutron star, it is thought to "bury" the magnetic field of the neutron star (although the details are unclear), leaving millisecond pulsars with magnetic fields 1000–10,000 times weaker than average pulsars. This low magnetic field is less effective at slowing the pulsar's rotation, so millisecond pulsars live for billions of years, making them the oldest known pulsars. Millisecond pulsars are seen in globular clusters, which stopped forming neutron stars billions of years ago.[28]
|
40 |
+
|
41 |
+
Of interest to the study of the state of the matter in a neutron
|
42 |
+
star are the glitches observed in the rotation velocity
|
43 |
+
of the neutron star. This velocity is decreasing slowly but steadily, except by sudden variations. One model put forward to explain these glitches is that they are the result of "starquakes" that adjust the crust of the neutron star. Models where the glitch is due to a decoupling of the possibly superconducting interior of the star have also been advanced. In both cases, the star's moment of inertia changes, but its angular momentum does not, resulting in a change in rotation rate.
|
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+
|
45 |
+
When two massive stars are born close together from the same cloud of gas, they can form a binary system and orbit each other from birth. If those two stars are at least a few times as massive as our sun, their lives will both end in supernova explosions. The more massive star explodes first, leaving behind a neutron star. If the explosion does not kick the second star away, the binary system survives. The neutron star can now be visible as a radio pulsar, and it slowly loses energy and spins down. Later, the second star can swell up, allowing the neutron star to suck up its matter. The matter falling onto the neutron star spins it up and reduces its magnetic field. This is called "recycling" because it returns the neutron star to a quickly-spinning state. Finally, the second star also explodes in a supernova, producing another neutron star. If this second explosion also fails to disrupt the binary, a double neutron star binary is formed. Otherwise, the spun-up neutron star is left with no companion and becomes a "disrupted recycled pulsar", spinning between a few and 50 times per second.[30]
|
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+
|
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+
The discovery of pulsars allowed astronomers to study an object never observed before, the neutron star. This kind of object is the only place where the behavior of matter at nuclear density can be observed (though not directly). Also, millisecond pulsars have allowed a test of general relativity in conditions of an intense gravitational field.
|
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|
49 |
+
Pulsar maps have been included on the two Pioneer plaques as well as the Voyager Golden Record. They show the position of the Sun, relative to 14 pulsars, which are identified by the unique timing of their electromagnetic pulses, so that our position both in space and in time can be calculated by potential extraterrestrial intelligences.[31] Because pulsars are emitting very regular pulses of radio waves, its radio transmissions do not require daily corrections. Moreover, pulsar positioning could create a spacecraft navigation system independently, or be used in conjunction with satellite navigation.[32][33]
|
50 |
+
|
51 |
+
Generally, the regularity of pulsar emission does not rival the stability of atomic clocks.[34] They can still be used as external reference.[35] For example, J0437-4715 has a period of 0.005757451936712637 s with an error of 1.7×10−17 s.
|
52 |
+
This stability allows millisecond pulsars to be used in establishing ephemeris time[36]
|
53 |
+
or in building pulsar clocks.[37]
|
54 |
+
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55 |
+
Timing noise is the name for rotational irregularities observed in all pulsars. This timing noise is observable as random wandering in the pulse frequency or phase.[38] It is unknown whether timing noise is related to pulsar glitches.
|
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+
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57 |
+
The radiation from pulsars passes through the interstellar medium (ISM) before reaching Earth. Free electrons in the warm (8000 K), ionized component of the ISM and H II regions affect the radiation in two primary ways. The resulting changes to the pulsar's radiation provide an important probe of the ISM itself.[39]
|
58 |
+
|
59 |
+
Because of the dispersive nature of the interstellar plasma, lower-frequency radio waves travel through the medium slower than higher-frequency radio waves. The resulting delay in the arrival of pulses at a range of frequencies is directly measurable as the dispersion measure of the pulsar. The dispersion measure is the total column density of free electrons between the observer and the pulsar,
|
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61 |
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where
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65 |
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D
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68 |
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{\displaystyle D}
|
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|
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+
is the distance from the pulsar to the observer and
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n
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e
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|
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{\displaystyle n_{e}}
|
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84 |
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is the electron density of the ISM. The dispersion measure is used to construct models of the free electron distribution in the Milky Way.[40]
|
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86 |
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Additionally, turbulence in the interstellar gas causes density inhomogeneities in the ISM which cause scattering of the radio waves from the pulsar. The resulting scintillation of the radio waves—the same effect as the twinkling of a star in visible light due to density variations in the Earth's atmosphere—can be used to reconstruct information about the small scale variations in the ISM.[41] Due to the high velocity (up to several hundred km/s) of many pulsars, a single pulsar scans the ISM rapidly, which results in changing scintillation patterns over timescales of a few minutes.[42]
|
87 |
+
|
88 |
+
Pulsars orbiting within the curved space-time around Sgr A*, the supermassive black hole at the center of the Milky Way, could serve as probes of gravity in the strong-field regime.[43] Arrival times of the pulses would be affected by special- and general-relativistic Doppler shifts and by the complicated paths that the radio waves would travel through the strongly curved space-time around the black hole. In order for the effects of general relativity to be measurable with current instruments, pulsars with orbital periods less than about 10 years would need to be discovered;[43] such pulsars would orbit at distances inside 0.01 pc from Sgr A*. Searches are currently underway; at present, five pulsars are known to lie within 100 pc from Sgr A*.[44]
|
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90 |
+
There are 3 consortia around the world which use pulsars to search for gravitational waves. In Europe, there is the European Pulsar Timing Array (EPTA); there is the Parkes Pulsar Timing Array (PPTA) in Australia; and there is the North American Nanohertz Observatory for Gravitational Waves (NANOGrav) in Canada and the US. Together, the consortia form the International Pulsar Timing Array (IPTA). The pulses from Millisecond Pulsars (MSPs) are used as a system of Galactic clocks. Disturbances in the clocks will be measurable at Earth. A disturbance from a passing gravitational wave will have a particular signature across the ensemble of pulsars, and will be thus detected.
|
91 |
+
|
92 |
+
The pulsars listed here were either the first discovered of its type, or represent an extreme of some type among the known pulsar population, such as having the shortest measured period.
|
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+
|
94 |
+
Video – Crab Pulsar – bright pulse & interpulse.
|
95 |
+
|
96 |
+
Video – Vela pulsar – X-ray light.
|
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+
|
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+
Video – Artist's impression of AR Scorpii.
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Pyongyang (US: /ˌpjɒŋˈjæŋ/, UK: /ˌpjʌŋˈjɑːŋ/,[6] Korean: [pʰjʌŋ.jaŋ]) is the capital and largest city of North Korea. Pyongyang is located on the Taedong River about 109 kilometers (68 mi) upstream from its mouth on the Yellow Sea. According to the 2008 population census, it has a population of 3,255,288.[7] Pyongyang is a directly-administered city (직할시; 直轄市; chikhalsi) with equal status to North Korean provinces.
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Pyongyang is considered one of the oldest cities in Korea.[8] It was the capital of two ancient Korean kingdoms, including Gojoseon and Goguryeo, and served as the secondary capital of Goryeo. Much of the city was destroyed during the First Sino-Japanese War, but it was revived under Japanese rule and became an industrial center. Following the establishment of North Korea in 1948, Pyongyang became its de facto capital. The city was again devastated during the Korean War, but was quickly rebuilt after the war with Soviet assistance.
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Pyongyang is the political, industrial and transport center of North Korea. It is home to North Korea's major government institutions, as well as the ruling Workers' Party of Korea.
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The city's other historic names include Kisong, Hwangsong, Rakrang, Sŏgyong, Sodo, Hogyong, Changan,[citation needed] and Heijō[9][10] (during Japanese rule in Korea). There are several variants.[a] During the early 20th century, Pyongyang came to be known among missionaries as being the "Jerusalem of the East", due to its historical status as a stronghold of Christianity, namely Protestantism, especially during the Pyongyang revival of 1907.[22][23]
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After Kim Il-sung's death in 1994, some members of Kim Jong-il's faction proposed changing the name of Pyongyang to "Kim Il-sung City" (Korean: 김일성시; Hanja: 金日成市), but others suggested that North Korea should begin calling Seoul "Kim Il-sung City" instead and grant Pyongyang the moniker "Kim Jong-il City", and in the end neither proposal was implemented.[24]
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In 1955, archaeologists excavated evidence of prehistoric occupation in a large ancient village in the Pyongyang area, called Kŭmtan-ni, dating to the Jeulmun and Mumun pottery periods.[25] North Koreans associate Pyongyang with the mythological city of "Asadal" (Korean: 아사달; Hanja: 阿斯達), or Wanggeom-seong (Korean: 왕검성; Hanja: 王儉城), the first second millennium BC capital of Gojoseon ("Old Joseon") according to Korean historiographies beginning with the 13th-century Samgungnyusa.
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Historians[who?] deny this claim because earlier Chinese historiographical works such as the Guanzi, Classic of Mountains and Seas, Records of the Grand Historian, and Records of the Three Kingdoms, mention a much later "Joseon".[citation needed] The connection between the two therefore may have been asserted by North Korea for the use of propaganda.[citation needed] Nevertheless, Pyongyang became a major city in old Joseon.
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Korean mythology asserts that Pyongyang was founded in 1122 BC on the site of the capital of the legendary king Dangun.[8] Wanggeom-seong, which was in the location of Pyongyang, became the capital of Gojoseon from 194 to 108 BC. It fell in the Han conquest of Gojoseon in 108 BC. Emperor Wu of Han ordered four commanderies be set up, with Lelang Commandery in the center and its capital established as "樂浪" (Old Chinese: *[r]ˤawk*[r]ˤaŋ,[26] Standard Chinese: pinyin: Lèlàng, Korean: 낙랑; Hanja: 樂浪). Several archaeological findings from the later, Eastern Han (20–220 AD) period in the Pyeongyang area seems to suggest that Han forces later launched brief incursions around these parts.
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The area around the city was called Nanglang during the early Three Kingdoms period. As the capital of Nanglang (Korean: 낙랑국; Hanja: 樂浪國),[b] Pyeongyang remained an important commercial and cultural outpost after the Lelang Commandery was destroyed by an expanding Goguryeo in 313.
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Goguryeo moved its capital there in 427. According to Christopher Beckwith, Pyongyang is the Sino-Korean reading of the name they gave it in their language: Piarna, or "level land".[27]
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In 668, Pyongyang became the capital of the Protectorate General to Pacify the East established by the Tang dynasty of China. However, by 676, it was taken by Silla, but left on the border between Silla and Balhae. Pyongyang was left abandoned during the Later Silla period, until it was recovered by Wang Geon and decreed as the Western Capital of Goryeo. During the Joseon period, it became the provincial capital of Pyeongan Province.
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During the Japanese invasions of Korea (1592–98), Pyongyang was captured by the Japanese and held until they were defeated in the Siege of Pyongyang.[8] Later in the 17th century, it became temporarily occupied during the Qing invasion of Joseon until peace arrangements were made between Korea and Qing China. While the invasions made Koreans suspicious of foreigners, the influence of Christianity began to grow after the country opened itself up to foreigners in the 16th century. Pyongyang became the base of Christian expansion in Korea, and by 1880 it had more than 100 churches and more Protestant missionaries than any other Asian city.[8]
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In 1890, the city had 40,000 inhabitants.[28] It was the site of the Battle of Pyongyang during the First Sino-Japanese War, which led to the destruction and depopulation of much of the city. It was the provincial capital of South Pyeongan Province beginning in 1896. Under Japanese colonial rule, the city became an industrial center, called Heijō (with the same Chinese characters 平壤 but read as へいじょう) in Japanese.
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In July 1931, the city experienced anti-Chinese riots as a result of the Wanpaoshan Incident and the sensationalized media reports about it which appeared in Imperial Japanese and Korean newspapers.[29]
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By 1938, Pyongyang had a population of 235,000.[28]
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On 25 August 1945, the Soviet 25th Army entered Pyongyang and it became the temporary capital of the Provisional People's Committee for North Korea. A People's Committee was already established there, led by veteran Christian nationalist Cho Man-sik.[30] Pyongyang became the de facto capital of North Korea upon its establishment in 1948. At the time, the Pyongyang government aimed to recapture Korea's official capital, Seoul. Pyongyang was again severely damaged in the Korean War, during which it was briefly occupied by South Korean forces from 19 October to 6 December 1950. In 1952, it was the target of the largest aerial raid of the entire war, involving 1,400 UN aircraft.
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Already during the war, plans were made to reconstruct the city. On 27 July 1953 – the day the armistice between North Korea and South Korea was signed – The Pyongyang Review wrote: "While streets were in flames, an exhibition showing the general plan of restoration of Pyongyang was held at the Moranbong Underground Theater", the air raid shelter of the government under Moranbong. "On the way of victory... fireworks which streamed high into the night sky of the capital in a gun salute briefly illuminated the construction plan of the city which would rise soon with a new look".[31] After the war, the city was quickly rebuilt with assistance from the Soviet Union, and many buildings were built in the style of Stalinist architecture. The plans for the modern city of Pyongyang were first displayed for public viewing in a theatre building. Kim Jung-hee, one of the founding members of the Korean Architects Alliance, who has studied architecture in prewar Japan was appointed by Kim Il Sung to build the city masterplan, Moscow Architecture Institute, designed the “Pyongyang City Reconstruction and Construction Comprehensive Plan” in 1951, and was officially adopted in 1953. The transformation into a modern, propaganda-designed city called Stalin style architecture with a Korean-style arrangement (and other modernist architecture that was said to have been greatly influenced by Brazilian architect Oscar Niemeyer) began.[32] and in the 1972 Constitution it was officially declared as the capital.
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In 2001, North Korean authorities began a long-term modernisation programme. The Ministry of Capital City Construction Development was included in the Cabinet in that year. In 2006, Kim Jong-il's brother-in-law Jang Song-thaek took charge of the ministry.
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Pyongyang is in the west-central part of North Korea; the city lies on a flat plain about 50 kilometres (31 mi) east of the Korea Bay, an arm of the Yellow Sea. The Taedong River flows southwestward through the city toward the Korea Bay. The Pyongyang plain, where the city is situated, is one of the two large plains on the Western coast of the Korean peninsula, the other being the Chaeryong plain. Both have an area of approximately 500 square kilometers.[33]
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Pyongyang has a hot-summer humid continental climate (Köppen: Dwa), featuring warm to hot, humid summers and cold, dry winters.[34][35] Cold, dry winds can blow from Siberia in winter, making conditions very cold; the low temperature is usually below freezing between November and early March, although the average daytime high is at least a few degrees above freezing in every month except January. The winter is generally much drier than summer, with snow falling for 37 days on average.
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The transition from the cold, dry winter to the warm, wet summer occurs rather quickly between April and early May, and there is a similarly abrupt return to winter conditions in late October and November. Summers are generally hot and humid, with the East Asian monsoon taking place from June until September; these are also the hottest months, with average temperatures of 21 to 25 °C (70 to 77 °F), and daytime highs often above 30 °C (86 °F).
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Major government and other public offices are located in Pyongyang, which is constitutionally designated as the country's capital.[39] The seat of the Workers' Party Central Committee and the Pyongyang People's Committee are located in Haebangsan-dong, Chung-guyok. The Cabinet of North Korea is located in Jongro-dong, Chung-guyok.
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Pyongyang is also the seat of all major North Korean security institutions. The largest of them, the Ministry of People's Security, has 130,000 employees working in 12 bureaus. These oversee activities including: police services, security of party officials, classified documents, census, civil registrations, large-scale public construction, traffic control, fire safety, civil defense, public health and customs.[40] Another significant structure based in the city is the State Security Department, whose 30,000 personnel manage intelligence, political prison systems, military industrial security and entry and exit management.[41]
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The politics and management of the city is dominated by the Workers' Party of Korea, as they are in the national level. The city is managed by the Pyongyang Party Committee of the Workers' Party of Korea and its chairman is the de facto mayor. The supreme standing state organ is the Pyongyang People's Committee, responsible
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for everyday events in support of the city. This includes following local Party guidance as channeled through the Pyongyang Party Committee, the distribution of resources prioritised to Pyongyang, and providing support to KWP and internal security agency personnel and families.
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P'yŏngyang is divided into 18 wards (ku- or guyŏk) (the city proper) and 2 counties (kun or gun).[42]
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Foreign media reports in 2010 stated that Kangnam-gun, Chunghwa-gun, Sangwŏn-gun, and Sŭngho-guyŏk had been transferred to the administration of neighboring North Hwanghae province.[43] However, Kangnam-gun was returned to Pyongyang in 2011.[44]
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After being destroyed during the Korean War, Pyongyang was entirely rebuilt according to Kim Il-sung's vision, which was to create a capital that would boost morale in the post-war years.[45] The result was a city with wide, tree-lined boulevards and public buildings with terraced landscaping, mosaics and decorated ceilings.[46] Its Russian-style architecture makes it reminiscent of a Siberian city during winter snowfall, although edifices of traditional Korean design somewhat soften this perception. In summer, it is notable for its rivers, willow trees, flowers and parkland.[46]
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The streets are laid out in a north-south, east-west grid, giving the city an orderly appearance.[46] North Korean designers applied the Swedish experience of self-sufficient urban neighbourhoods throughout the entire country, and Pyongyang is no exception. Its inhabitants are mostly divided into administrative units of 5,000 to 6,000 people (dong). These units all have similar sets of amenities including a food store, a barber shop, a tailor, a public bathhouse, a post office, a clinic, a library and others. Many residents occupy high-rise apartment buildings.[47] One of Kim Il-sung's priorities while designing Pyongyang was to limit the population. Authorities maintain a restrictive regime of movement into the city, making it atypical of East Asia as it is silent, uncrowded and spacious.[48]
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Structures in Pyongyang are divided into three major architectural categories: monuments, buildings with traditional Korean motifs and high-rises.[49] Some of North Korea's most recognisable landmarks are monuments, like the Juche Tower, the Arch of Triumph and the Mansu Hill Grand Monument. The first of them is a 170-meter granite spire symbolizing the Juche ideology. It was completed in 1982 and contains 25,550 granite blocks, one for each day of Kim Il-sung's life up to that point.[49] The most prominent building on Pyongyang's skyline is Ryugyong Hotel,[49] the seventh highest building in the world terms of floor count, the tallest unoccupied building in the world,[50] and one of the tallest hotels in the world. It has yet to open.[51][52]
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Pyongyang has a rapidly evolving skyline, dominated by high-rise apartment buildings. A construction boom began with the Changjon Street Apartment Complex, which was completed in 2012.[53] Construction of the complex began after late leader Kim Jong-il described Changjon Street as "pitiful".[54] Other housing complexes are being upgraded as well, but most are still poorly insulated, and lacking elevators and central heating.[55] An urban renewal program continued under Kim Jong-un's leadership, with the old apartments of the 1970s and '80s replaced by taller high rise buildings and leisure parks like the Kaesong Youth Park, as well as renovations of older buildings.[56] In 2018, the city was described as unrecognizable compared to five years before.[57]
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Notable landmarks in the city include:
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Pyongyang TV Tower is a minor landmark. Other visitor attractions include the Korea Central Zoo. The Arch of Reunification has a map of a united Korea supported by two concrete Korean women dressed in traditional dress straddling the Reunification Highway, which stretches from Pyongyang to the Korean Demilitarized Zone (DMZ).
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Juche Tower Monument to the philosophy of Juche (self-reliance)
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Arch of Triumph
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Arch of Reunification, a monument to the goal of a reunified Korea
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Monument to Party Founding
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Rungrado May Day Stadium
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Kumsusan Palace of the Sun
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Tomb of King Tongmyeong
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Ryugyong Hotel
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Mansu Hill Grand Monument
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Mirae (Future) Scientist's Street
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Pyongyang served as the provincial capital of South Pyongan Province until 1946,[58] and Pyongyang cuisine shares the general culinary tradition of the Pyongan province. The most famous local food is Pyongyang raengmyŏn, or also called mul raengmyŏn or just simply raengmyŏn. Raengmyŏn literally means "cold noodles", while the affix mul refers to water because the dish is served in a cold broth. Raengmyŏn consists of thin and chewy buckwheat noodles in a cold meat-broth with dongchimi (watery kimchi) and topped with a slice of sweet Korean pear.
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Pyongyang raengmyŏn was originally eaten in homes built with ondol (traditional underfloor heating) during the cold winter, so it is also called "Pyongyang deoldeori" (shivering in Pyongyang). Pyongyang locals sometimes enjoyed it as a haejangguk, which is any type of food eaten as a hangover-cure, usually a warm soup.[59]
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Another representative Pyongyang dish, Taedonggang sungeoguk, translates as "trout soup from the Taedong River". The soup features trout (abundant in the Taedong River) along with black peppercorns and salt.[60] Traditionally, it has been served to guests visiting Pyongyang. Therefore, there is a common saying, "How good was the trout soup?", which is used to greet people returning from Pyongyang. Another local specialty, Pyongyang onban (literally "warm rice of Pyongyang") comprises freshly cooked rice topped with sliced mushrooms, chicken, and a couple of bindaetteok (pancakes made from ground mung beans and vegetables).[59]
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In 2018, there were many high quality restaurants in Pyongyang with Korean and international food, and imported alcoholic beverages.[61] Famous restaurants include Okryu-gwan and Ch'ongryugwan.[62] Some street foods exist in Pyongyang, where vendors operate food stalls.[63] Foreign foods like hamburgers, fries, pizza, and coffee are easily found.[61] There is an active nightlife with late-night restaurants and karaoke.[61]
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The city has water parks, amusement parks, skating rinks, health clubs, a shooting range, and a dolphinarium.[56]
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Pyongyang has a number of sports clubs, including the April 25 Sports Club and the Pyongyang City Sports Club.[64] The most popular sport in Pyongyang is football.[citation needed]
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Pyongyang is North Korea's industrial center.[8] Thanks to the abundance of natural resources like coal, iron and limestone, as well as good land and water transport systems, it was the first industrial city to emerge in North Korea after the Korean War. Light and heavy industries are both present and have developed in parallel. Heavy manufactures include cement, industrial ceramics, munitions and weapons, but mechanical engineering remains the core industry. Light industries in Pyongyang and its vicinity include textiles, footwear and food, among others. Special emphasis is put on the production and supply of fresh produce and subsidiary crops in farms on the city's outskirts. Other crops include rice, sweetcorn and soybeans. Pyongyang aims to achieve self-sufficiency in meat production. High-density facilities raise pigs, chicken and other livestock.[8]
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The city still experiences frequent shortages of electricity.[65] To solve this problem, two power stations – Huichon Power Stations 1 and 2 – were built in Chagang Province and supply the city through direct transmission lines. A second phase of the power expansion project was launched in January 2013, consisting of a series of small dams along the Chongchon River. The first two power stations have a maximum generating capacity of 300 megawatts (MW), while the 10 dams to be built under second phase are expected to generate about 120 MW.[65] In addition, the city has several existing or planned thermal power stations. These include Pyongyang TPS with a capacity of 500 MW, East Pyongyang TPS with a capacity of 50 MW, and Kangdong TPS which is under construction.[66]
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Pyongyang is home to several large department stores including the Pothonggang Department Store, Pyongyang Department Store No. 1, Pyongyang Department Store No. 2, Kwangbok Department Store, Ragwon Department Store, Pyongyang Station Department Store, and the Pyongyang Children's Department Store.[67]
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The city also has Hwanggumbol Shop, a chain of state-owned convenience stores supplying goods at prices cheaper than those in the jangmadang markets. Hwanggumbol Shops are specifically designed to control North Korea's expanding markets by attracting consumers and guaranteeing the circulation of money in government-operated stores.[68]
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Pyongyang is also the main transport hub of the country: it has a network of roads, railways and air routes which link it to both foreign and domestic destinations. It is the starting point of inter-regional highways reaching Nampo, Wonsan and Kaesong.[8] Pyongyang railway station serves the main railway lines, including the Pyongui Line and the Pyongbu Line. Regular international rail services to Beijing, the Chinese border city of Dandong and Moscow are also available.
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A rail journey to Beijing takes about 25 hours and 25 minutes (K27 from Beijing/K28 from Pyongyang, on Mondays, Wednesdays, Thursdays and Saturdays); a journey to Dandong takes about 6 hours (daily); a journey to Moscow takes six days. The city also connects to the Eurasian Land Bridge via the Trans-Siberian Railway. A high-speed rail link to Wonsan is planned.[69]
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The Metro, tram and trolleybus systems are used mainly by commuters as a primary means of urban transportation.[8] Cycle lanes were introduced on main thoroughfares in July 2015.[70] There are relatively few cars in the city. Cars are a symbol of status in the country due to their scarcity as a result of restrictions on import because of international sanctions and domestic regulations.[71] Some roads are also reported to be in poor condition.[72] However, by 2018, Pyongyang had begun to experience traffic jams.[61]
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State-owned Air Koryo has scheduled international flights from Pyongyang Sunan International Airport to Beijing (PEK), Shenyang (SHE), Vladivostok (VVO), Shanghai (PVG) and Dandong.[73] The only domestic destinations are Hamhung,
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Wonsan, Chongjin, Hyesan and Samjiyon. Since 31 March 2008, Air China launched a regular service between Beijing and Pyongyang,[74] although Air China's flights are often canceled due to lack of passengers.[75]
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Kim Il-sung University, North Korea's oldest university, was established in 1946.[8] It has seven colleges, 14 faculties and 16 other institutes, graduate schools and university units.[76] These include the primary medical education and health personnel training unit, the medical college; a physics faculty which covers a range of studies including theoretical physics, optical science, geophysics and astrophysics;[77] an atomic energy institute and a human evolution research office which studies human evolution through a Juche point of view. Kim Il-sung University also has its own publishing house, sports club (Ryongnamsan Sports Club),[78] revolutionary museum, nature museum, libraries, a gym, indoor swimming pool and educator apartment houses. Its two main buildings were completed in 1965 (Building 1) and 1972 (Building 2). A third building on campus is planned.[79]
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Other higher education establishments include Kim Chaek University of Technology, Pyongyang University of Music and Dance and Pyongyang University of Foreign Studies. Pyongyang University of Science and Technology (PUST) is the country's first private university where most of the lecturers are American and courses are carried out in English.[80][81] A science and technology hall is under construction on Ssuk Islet. Its stated purpose is to contribute to the "informatization of educational resources" by centralizing teaching materials, compulsory literature and experimental data for state-level use in a digital format.[82]
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Sosong-guyok hosts a 20 MeV cyclotron called MGC-20. The initial project was approved by the International Atomic Energy Agency (IAEA) in 1983 and funded by the IAEA, the United States and the North Korean government. The cyclotron was ordered from the Soviet Union in 1985 and constructed between 1987 and 1990. It is used for student training, production of medical isotopes for nuclear medicine as well as studies in biology, chemistry and physics.[83]
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Medical centers include the Red Cross Hospital, the First People's Hospital which is located near Moran Hill and was the first hospital to be built in North Korea after the liberation of Korea in 1945,[84] the Second People's Hospital, Ponghwa Recuperative Center (also known as Bonghwa Clinic or Presidential Clinic) located in Sokam-dong, Potonggang-guyok, 1.5 km (0.93 mi) northwest of Kim Il-sung Square,[85] Pyongyang Medical School Hospital, Namsan Treatment Center which is adjacent[86] Pyongyang's Maternity Hospital, Taesongsan General Hospital,[87] Kim Man-yoo Hospital, Staff Treatment Center and Okryu Children's Hospital. A new hospital named Pyongyang General Hospital is under-construction in Pyongyang.[88]
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Pyongyang is twinned with:
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The bald eagle (Haliaeetus leucocephalus) is a bird of prey found in North America. A sea eagle, it has two known subspecies and forms a species pair with the white-tailed eagle (Haliaeetus albicilla). Its range includes most of Canada and Alaska, all of the contiguous United States, and northern Mexico. It is found near large bodies of open water with an abundant food supply and old-growth trees for nesting.
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The bald eagle is an opportunistic feeder which subsists mainly on fish, which it swoops down and snatches from the water with its talons. It builds the largest nest of any North American bird and the largest tree nests ever recorded for any animal species, up to 4 m (13 ft) deep, 2.5 m (8.2 ft) wide, and 1 metric ton (1.1 short tons) in weight. Sexual maturity is attained at the age of four to five years.
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Bald eagles are not actually bald; the name derives from an older meaning of the word, "white headed". The adult is mainly brown with a white head and tail. The sexes are identical in plumage, but females are about 25 percent larger than males. The beak is large and hooked. The plumage of the immature is brown.
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The bald eagle is the national bird of the United States of America. The bald eagle appears on its seal. In the late 20th century it was on the brink of extirpation in the contiguous United States. Populations have since recovered and the species was removed from the U.S. government's list of endangered species on July 12, 1995 and transferred to the list of threatened species. It was removed from the List of Endangered and Threatened Wildlife in the contiguous states on June 28, 2007.
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The plumage of an adult bald eagle is evenly dark brown with a white head and tail. The tail is moderately long and slightly wedge-shaped. Males and females are identical in plumage coloration, but sexual dimorphism is evident in the species, in that females are 25% larger than males.[3] The beak, feet and irises are bright yellow. The legs are feather-free, and the toes are short and powerful with large talons. The highly developed talon of the hind toe is used to pierce the vital areas of prey while it is held immobile by the front toes.[4] The beak is large and hooked, with a yellow cere.[5] The adult bald eagle is unmistakable in its native range. The closely related African fish eagle (Haliaeetus vocifer) (from far outside the bald eagle's range) also has a brown body (albeit of somewhat more rufous hue), white head and tail, but differs from the bald in having a white chest and black tip to the bill.[6]
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The plumage of the immature is a dark brown overlaid with messy white streaking until the fifth (rarely fourth, very rarely third) year, when it reaches sexual maturity.[3][4] Immature bald eagles are distinguishable from the golden eagle (Aquila chrysaetos), the only other very large, non-vulturine raptorial bird in North America, in that the former has a larger, more protruding head with a larger beak, straighter edged wings which are held flat (not slightly raised) and with a stiffer wing beat and feathers which do not completely cover the legs. When seen well, the golden eagle is distinctive in plumage with a more solid warm brown color than an immature bald eagle, with a reddish-golden patch to its nape and (in immature birds) a highly contrasting set of white squares on the wing.[7] Another distinguishing feature of the immature bald eagle over the mature bird is its black, yellow-tipped beak; the mature eagle has a fully yellow beak. The bald eagle has sometimes been considered the largest true raptor (accipitrid) in North America. The only larger species of raptor-like bird is the California condor (Gymnogyps californianus), a New World vulture which today is not generally considered a taxonomic ally of true accipitrids.[8] However, the golden eagle, averaging 4.18 kg (9.2 lb) and 63 cm (25 in) in wing chord length in its American race (Aquila chrysaetos canadensis), is merely 455 g (1.003 lb) lighter in mean body mass and exceeds the bald eagle in mean wing chord length by around 3 cm (1.2 in).[6][9] Additionally, the bald eagle's close cousins, the relatively longer-winged but shorter-tailed white-tailed eagle and the overall larger Steller's sea eagle (Haliaeetus pelagicus), may, rarely, wander to coastal Alaska from Asia.[6]
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The bald eagle has a body length of 70–102 cm (28–40 in). Typical wingspan is between 1.8 and 2.3 m (5 ft 11 in and 7 ft 7 in) and mass is normally between 3 and 6.3 kg (6.6 and 13.9 lb).[6] Females are about 25% larger than males, averaging as much as 5.6 kg (12 lb), and against the males' average weight of 4.1 kg (9.0 lb).[3][10][11][12]
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The size of the bird varies by location and generally corresponds with Bergmann's rule, since the species increases in size further away from the Equator and the tropics. For example, eagles from South Carolina average 3.27 kg (7.2 lb) in mass and 1.88 m (6 ft 2 in) in wingspan, smaller than their northern counterparts.[13] One field guide in Florida listed similarly small sizes for bald eagles there, at about 4.13 kg (9.1 lb).[14] Of intermediate size, 117 migrant bald eagles in Glacier National Park were found to average 4.22 kg (9.3 lb) but this was mostly (possibly post-dispersal) juvenile eagles, with 6 adults here averaging 4.3 kg (9.5 lb).[15] Wintering eagles in Arizona (winter weights are usually the highest through the year since like many raptors they spend the highest percentage of time foraging during winter) were found to average 4.74 kg (10.4 lb).[16] The largest eagles are from Alaska, where large females may weigh more than 7 kg (15 lb) and span 2.44 m (8 ft 0 in) across the wings.[5][17] A survey of adult weights in Alaska showed that females there weighed on average 5.35 kg (11.8 lb), respectively, and males weighed 4.23 kg (9.3 lb) against immatures which averaged 5.09 kg (11.2 lb) and 4.05 kg (8.9 lb) in the two sexes.[18][19] An Alaskan adult female eagle that was considered outsized weighed some 7.4 kg (16 lb).[20] R.S. Palmer listed a record from 1876 in Wyoming County, New York of an enormous adult bald eagle that was shot and reportedly scaled 8.2 kg (18 lb).[19] Among standard linear measurements, the wing chord is 51.5–69 cm (20.3–27.2 in), the tail is 23–37 cm (9.1–14.6 in) long, and the tarsus is 8 to 11 cm (3.1 to 4.3 in).[6][21] The culmen reportedly ranges from 3 to 7.5 cm (1.2 to 3.0 in), while the measurement from the gape to the tip of the bill is 7–9 cm (2.8–3.5 in).[21][22] The bill size is unusually variable as Alaskan eagles could be up to twice the bill length of "southern birds" (i.e. from Georgia, Louisiana, Florida), with means in between the sexes of 6.83 cm (2.69 in) and 4.12 cm (1.62 in) in culmen length, respectively, from these two areas.[23][24]
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The call consists of weak staccato, chirping whistles, kleek kik ik ik ik, somewhat similar in cadence to a gull's call. The calls of young birds tend to be more harsh and shrill than those of adults.[6][7]
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The bald eagle is placed in the genus Haliaeetus (sea eagles), and gets both its common and specific scientific names from the distinctive appearance of the adult's head. Bald in the English name is from the older usage meaning "white" rather than "hairless", referring to the white head and tail feathers and their contrast with the darker body, as in piebald.[25] The genus name is New Latin: Haliaeetus (from the Ancient Greek: ἁλιάετος, romanized: haliaetos, lit. 'sea eagle'),[26] and the specific name, leucocephalus, is Latinized (Ancient Greek: λευκός, romanized: leukos, lit. 'white')[27] and (κεφαλή, kephalḗ, 'head').[28][29]
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The bald eagle was one of the many species originally described by Carl Linnaeus in his 18th-century work Systema Naturae, under the name Falco leucocephalus.[30]
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There are two recognized subspecies of bald eagle:[3][31]
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The bald eagle forms a species pair with the white-tailed eagle of Eurasia. This species pair consists of a white-headed and a tan-headed species of roughly equal size; the white-tailed eagle also has overall somewhat paler brown body plumage. The two species fill the same ecological niche in their respective ranges. The pair diverged from other sea eagles at the beginning of the Early Miocene (c. 10 Ma BP) at the latest, but possibly as early as the Early/Middle Oligocene, 28 Ma BP, if the most ancient fossil record is correctly assigned to this genus.[33]
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The bald eagle's natural range covers most of North America, including most of Canada, all of the continental United States, and northern Mexico. It is the only sea eagle endemic to North America. Occupying varied habitats from the bayous of Louisiana to the Sonoran Desert and the eastern deciduous forests of Quebec and New England, northern birds are migratory, while southern birds are resident, remaining on their breeding territory all year. At minimum population, in the 1950s, it was largely restricted to Alaska, the Aleutian Islands, northern and eastern Canada, and Florida.[34] During the interval 1966-2015 bald eagle numbers increased substantially throughout its winter and breeding ranges,[35] and as of 2018 the species nests in every continental state and province in the United States and Canada.[36]
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The majority of bald eagles in Canada are found along the British Columbia coast while large populations are found in the forests of Alberta, Saskatchewan, Manitoba and Ontario.[37] Bald eagles also congregate in certain locations in winter. From November until February, one to two thousand birds winter in Squamish, British Columbia, about halfway between Vancouver and Whistler. The birds primarily gather along the Squamish and Cheakamus Rivers, attracted by the salmon spawning in the area.[38] Similar congregations of wintering bald eagles at open lakes and rivers, wherein fish are readily available for hunting or scavenging, are observed in the northern United States.[39]
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It has occurred as a vagrant twice in Ireland; a juvenile was shot illegally in Fermanagh on January 11, 1973 (misidentified at first as a white-tailed eagle), and an exhausted juvenile was captured in Kerry on November 15, 1987.[40]
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The bald eagle occurs during its breeding season in virtually any kind of American wetland habitat such as seacoasts, rivers, large lakes or marshes or other large bodies of open water with an abundance of fish. Studies have shown a preference for bodies of water with a circumference greater than 11 km (7 mi), and lakes with an area greater than 10 km2 (4 sq mi) are optimal for breeding bald eagles.[41]
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The bald eagle typically requires old-growth and mature stands of coniferous or hardwood trees for perching, roosting, and nesting. Tree species reportedly is less important to the eagle pair than the tree's height, composition and location.[42] Perhaps of paramount importance for this species is an abundance of comparatively large trees surrounding the body of water. Selected trees must have good visibility, be over 20 m (66 ft) tall, an open structure, and proximity to prey. If nesting trees are in standing water such as in a mangrove swamp, the nest can be located fairly low, at as low 6 m (20 ft) above the ground.[43] In a more typical tree standing on dry ground, nests may be located from 16 to 38 m (52 to 125 ft) in height. In Chesapeake Bay, nesting trees averaged 82 cm (32 in) in diameter and 28 m (92 ft) in total height, while in Florida, the average nesting tree stands 23 m (75 ft) high and is 23 cm (9.1 in) in diameter.[44][45] Trees used for nesting in the Greater Yellowstone area average 27 m (89 ft) high.[46] Trees or forest used for nesting should have a canopy cover of no more than 60%, and no less than 20%, and be in close proximity to water.[41] Most nests have been found within 200 m (660 ft) of open water. The greatest distance from open water recorded for a bald eagle nest was over 3 km (1.9 mi), in Florida.[8]
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Bald eagle nests are often very large in order to compensate for size of the birds. The largest recorded nest was found in Florida in 1963, and was measured at nearly 10 feet wide and 20 feet deep.[47]
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In Florida, nesting habitats often consist of mangrove swamps, the shorelines of lakes and rivers, pinelands, seasonally flooded flatwoods, hardwood swamps, and open prairies and pastureland with scattered tall trees. Favored nesting trees in Florida are slash pines (Pinus elliottii), longleaf pines (P. palustris), loblolly pines (P. taeda) and cypress trees, but for the southern coastal areas where mangroves are usually used.[43] In Wyoming, groves of mature cottonwoods or tall pines found along streams and rivers are typical bald eagle nesting habitats. Wyoming eagles may inhabit habitat types ranging from large, old-growth stands of ponderosa pines (Pinus ponderosa) to narrow strips of riparian trees surrounded by rangeland.[8] In Southeast Alaska, Sitka spruce (Picea sitchensis) provided 78% of the nesting trees used by eagles, followed by hemlocks (Tsuga) at 20%.[42] Increasingly, eagles nest in man-made reservoirs stocked with fish.[43]
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The bald eagle is usually quite sensitive to human activity while nesting, and is found most commonly in areas with minimal human disturbance. It chooses sites more than 1.2 km (0.75 mi) from low-density human disturbance and more than 1.8 km (1.1 mi) from medium- to high-density human disturbance.[41] However, bald eagles will occasionally nest in large estuaries or secluded groves within major cities, such as Hardtack Island on the Willamette River in Portland, Oregon or John Heinz National Wildlife Refuge at Tinicum in Philadelphia, Pennsylvania, which are surrounded by a great quantity of human activity.[48][49] Even more contrary to the usual sensitivity to disturbance, a family of bald eagles moved to the Harlem neighborhood in New York City in 2010.[50]
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While wintering, bald eagles tend to be less habitat and disturbance sensitive. They will commonly congregate at spots with plentiful perches and waters with plentiful prey and (in northern climes) partially unfrozen waters. Alternately, non-breeding or wintering bald eagles, particularly in areas with a lack of human disturbance, spend their time in various upland, terrestrial habitats sometimes quite far away from waterways. In the northern half of North America (especially the interior portion), this terrestrial inhabitance by bald eagles tends to be especially prevalent because unfrozen water may not be accessible. Upland wintering habitats often consist of open habitats with concentrations of medium-sized mammals, such as prairies, meadows or tundra, or open forests with regular carrion access.[8][42]
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The bald eagle is a powerful flier, and soars on thermal convection currents. It reaches speeds of 56–70 km/h (35–43 mph) when gliding and flapping, and about 48 km/h (30 mph) while carrying fish.[51] Its dive speed is between 120–160 km/h (75–99 mph), though it seldom dives vertically.[52] Regarding their flying abilities, despite being morphologically less well adapted to faster flight than golden eagles (especially during dives), the bald eagle is considered surprisingly maneuverable in flight. Bounty hunters shooting from helicopters opined that they were far more difficult to hunt while flying than golden eagles as they would turn, double back or dive as soon as approached. Bald eagles have also been recorded catching up to and then swooping under geese in flight, turning over and thrusting their talons into the other bird's breast.[19] It is partially migratory, depending on location. If its territory has access to open water, it remains there year-round, but if the body of water freezes during the winter, making it impossible to obtain food, it migrates to the south or to the coast. A number of populations are subject to post-breeding dispersal, mainly in juveniles; Florida eagles, for example, will disperse northwards in the summer.[53] The bald eagle selects migration routes which take advantage of thermals, updrafts, and food resources. During migration, it may ascend in a thermal and then glide down, or may ascend in updrafts created by the wind against a cliff or other terrain. Migration generally takes place during the daytime, usually between the local hours of 8:00 a.m. and 6:00 p.m., when thermals are produced by the sun.[4]
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The bald eagle is an opportunistic carnivore with the capacity to consume a great variety of prey. Throughout their range, fish often comprise the majority of the eagle's diet.[54] In 20 food habit studies across the species' range, fish comprised 56% of the diet of nesting eagles, birds 28%, mammals 14% and other prey 2%.[55] More than 400 species are known to be included in the bald eagle's prey spectrum, far more than its ecological equivalent in the Old World, the white-tailed eagle, is known to take. Despite its considerably lower population, the bald eagle may come in second amongst all North American accipitrids, slightly behind only the red-tailed hawk, in number of prey species recorded.[19][55][56][57] In Southeast Alaska, fish comprise approximately 66% of the year-around diet of bald eagles and 78% of the prey brought to the nest by the parents.[58] Eagles living in the Columbia River Estuary in Oregon were found to rely on fish for 90% of their dietary intake.[59] At least 100 species of fish have been recorded in the bald eagle's diet.[56] In the Pacific Northwest, spawning trout and salmon provide most of the bald eagles' diet from late summer throughout fall.[60] Southeast Alaskan eagles largely prey on pink salmon (Oncorhynchus gorbuscha), coho salmon (O. kisutch) and, more locally, sockeye salmon (O. nerka), with Chinook salmon (O. tshawytscha), due to their large size (12 to 18 kg (26 to 40 lb) average adult size) probably being taken only as carrion.[58] Also important in the estuaries and shallow coastlines of southern Alaska are Pacific herring (Clupea pallasii), Pacific sand lance (Ammodytes hexapterus) and eulachon (Thaleichthys pacificus).[58]
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In Oregon's Columbia River Estuary, the most significant prey species were largescale suckers (Catostomus macrocheilus) (17.3% of the prey selected there), American shad (Alosa sapidissima; 13%) and common carp (Cyprinus carpio; 10.8%).[59] Eagles living in the Chesapeake Bay in Maryland were found to subsist largely on American gizzard shad (Dorosoma cepedianum), threadfin shad (Dorosoma petenense) and white bass (Morone chrysops).[61] Floridian eagles have been reported to prey on catfish, most prevalently the brown bullhead (Ameiurus nebulosus) and any species in the genus Ictalurus as well as mullet, trout, needlefish, and eels.[8][43][62] Wintering eagles on the Platte River in Nebraska preyed mainly on American gizzard shads and common carp.[63] From observation in the Columbia River, 58% of the fish were caught alive by the eagle, 24% were scavenged as carcasses and 18% were pirated away from other animals.[59]
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Prey fish targeted by bald eagles are often quite large. When experimenters offered fish of different sizes in the breeding season around Lake Britton in California, fish measuring 34 to 38 cm (13 to 15 in) were taken 71.8% of the time by parent eagles while fish measuring 23 to 27.5 cm (9.1 to 10.8 in) were chosen only 25% of the time.[64] At nests around Lake Superior, the remains of fish (mostly suckers) were found to average 35.4 cm (13.9 in) in total length.[65] In the Columbia River estuary, most preyed on by eagles were estimated to measure between 30 and 60 cm (12 and 24 in) in length, and carp flown with (laboriously) were up to 86 cm (34 in) in length.[59] Much larger marine fish such as Pacific halibut (Hippoglossus stenolepis) and lemon sharks (Negaprion brevirostris) have been recorded among bald eagle prey though probably are only taken as young, as small, newly mature fish, or as carrion.[57][66]
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Benthic fishes such as catfish are usually consumed after they die and float to the surface, though while temporarily swimming in the open may be more vulnerable to predation than most fish since their eyes focus downwards.[61] Bald eagles also regularly exploit water turbines which produce battered, stunned or dead fish easily consumed.[67] Predators who leave behind scraps of dead fish that they kill, such as brown bears (Ursus arctos), gray wolves (Canis lupus) and red foxes (Vulpes vulpes), may be habitually followed in order to scavenge the kills secondarily.[58] Once North Pacific salmon die off after spawning, usually local bald eagles eat salmon carcasses almost exclusively. Eagles in Washington need to consume 489 g (1.078 lb) of fish each day for survival, with adults generally consuming more than juveniles and thus reducing potential energy deficiency and increasing survival during winter.[68]
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Behind fish, the next most significant prey base for bald eagles are other waterbirds. The contribution of such birds to the eagle's diet is variable, depending on the quantity and availability of fish near the water's surface. Waterbirds can seasonally comprise from 7% to 80% of the prey selection for eagles in certain localities.[59][69] Overall, birds are the most diverse group in the bald eagle's prey spectrum, with 200 prey species recorded.[19][56][57] Exceptionally, in the Greater Yellowstone area, birds were eaten as regularly as fish year-around, with both prey groups comprising 43% of the studied dietary intake.[46] Preferred avian prey includes grebes, alcids, ducks, gulls, coots, herons, egrets, and geese.[70]
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Bird species most preferred as prey by eagles tend to be medium-sized, such as western grebes (Aechmophorus occidentalis), mallards (Anas platyrhynchos) and American coots (Fulica americana) as such prey is relatively easy for the much larger eagles to catch and fly with.[8][59] American herring gull (Larus smithsonianus) are the favored avian prey species for eagles living around Lake Superior.[65] Larger waterbirds are occasionally prey as well, with wintering emperor geese (Chen canagica) and snow geese (C. caerulescens), which gather in large groups, sometimes becoming regular prey.[21][71] Other large waterbirds hunted at least occasionally by bald eagles have included adults of common loons (Gavis immer),[72] great black-backed gulls (Larus marinus),[73] sandhill cranes (Grus canadensis),[74] great blue herons (Ardea herodias),[55] Canada geese (Branta canadensis),[61] brown pelicans (Pelecanus occidentalis),[43] and fledgling American white pelicans (P. erythrorhynchos).[75] Colony nesting seabirds may be especially vulnerable to predation. Due to easy accessibility and lack of formidable nest defense by such species, bald eagles are capable of preying on such seabirds at all ages, from eggs to mature adults, and can effectively cull large portions of a colony.[76]
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Along some portions of the North Pacific coastline, bald eagles which had historically preyed mainly kelp-dwelling fish and supplementally sea otter (Enhydra lutris) pups are now preying mainly on seabird colonies since both the fish (possibly due to overfishing) and otters (cause unknown) have had precipitous population declines, causing concern for seabird conservation.[77] Because of this more extensive predation, some biologist have expressed concern that murres are heading for a "conservation collision" due to heavy eagle predation.[76] Eagles have been confirmed to attack nocturnally active, burrow-nesting seabird species such as storm petrels and shearwaters by digging out their burrows and feeding on all animals they find inside.[78] If a bald eagle flies close by, waterbirds will often fly away en masse, though in other cases they may seemingly ignore a perched eagle. If the said birds are on a colony, this exposed their unprotected eggs and nestlings to scavengers such as gulls.[76] Bird prey may occasionally be attacked in flight, with prey up to the size of Canada geese attacked and killed in mid-air.[70] Unprecedented photographs of a bald eagle unsuccessfully attempting to prey on a much larger adult trumpeter swan (Cygnus buccinator) in mid-flight were taken in 2012.[79] While adults often actively prey on waterbirds, congregated wintering waterfowl are frequently exploited for carcasses to scavenge by immature eagles in harsh winter weather.[80] Bald eagles have been recorded as killing other raptors on occasion. In some cases, these may be attacks of competition or kleptoparasitism on rival species but ended with the consumption of the victim. Nine species each of other accipitrids and owls are known to have been preyed upon by bald eagles. Owl prey species have ranged in size from western screech-owls (Megascops kennicotti) to snowy owls (Bubo scandiacus).[19][56][57][81] Larger diurnal raptors known to have fallen victim to bald eagles have included red-tailed hawks (Buteo jamaicensis),[82] peregrine falcons (Falco peregrinus),[83] northern goshawks (Accipiter gentilis),[84] ospreys (Pandion haliaetus)[85] and black (Coragyps atratus) and turkey vultures (Cathartes aura).[86]
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Mammalian prey includes rabbits, hares, ground squirrels, raccoons (Procyon lotor), muskrats (Ondatra zibethicus), beavers (Castor canadensis), and deer fawns. Newborn, dead, sickly or already injured mammals are often targeted. However, more formidable prey such as adult raccoons and subadult beavers are sometimes attacked. In the Chesapeake Bay area, bald eagles are reportedly the main natural predators of raccoons.[87][88] Other relatively large mammalian prey known to be taken by bald eagles (at least rarely) as adults include Virginia opossums (Didelphis virginiana), nine-banded armadillos (Dasypus novemcinctus), red and Arctic foxes (Vulpes vulpes & Vulpes lagopus) and striped skunks (Mephitis mephitis).[89][90] Even an adult bobcat (Lynx rufus) has been recorded amongst their prey, although this may have been scavenged.[91] Where available, seal colonies can provide much food. On Protection Island, Washington, they commonly feed on harbor seal (Phoca vitulina) afterbirths, still-borns and sickly seal pups.[92] On San Juan Island in Washington, introduced European rabbits (Oryctolagus cuniculus), mainly those killed by auto accidents, comprise nearly 60% of the dietary intake of eagles.[93] In landlocked areas of North America, wintering bald eagles may become habitual predators of medium-sized mammals that occur in colonies or local concentrations, such as prairie dogs (Cynomys) and jackrabbits (Lepus). Like the golden eagle, bald eagles are capable of attacking jackrabbits and hares of nearly any size[8][94] Together with the golden eagle, bald eagles are occasionally accused of preying on livestock, especially sheep (Ovis aries). There are a handful of proven cases of lamb predation, some of specimens weighing up to 11 kg (24 lb), by bald eagles but they are much less likely to attack a healthy lamb than a golden eagle and both species prefer native, wild prey and are unlikely to cause any extensive detriment to human livelihoods.[95] There is one case of a bald eagle killing and feeding on an adult, pregnant ewe (then joined in eating the kill by at least 3 other eagles), which, weighing on average over 60 kg (130 lb), is much larger than any other known prey taken by this species.[96]
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Supplemental prey are readily taken given the opportunity. In some areas reptiles may become regular prey, especially warm areas such as Florida where reptile diversity is high. Turtles are perhaps the most regularly hunted type of reptile.[8] In coastal New Jersey, 14 of 20 studied eagle nests included remains of turtles. The main species found were common musk turtles (Sternotherus odoratus), diamondback terrapin (Malaclemys terrapin) and juvenile common snapping turtles (Chelydra serpentina). In these New Jersey nests, mainly subadult and small adults were taken, ranging in carapace length from 9.2 to 17.1 cm (3.6 to 6.7 in).[97] Similarly, many turtles were recorded in the diet in the Chesapeake Bay.[98] Snakes are also taken occasionally, especially partially aquatic ones, as are amphibians and crustaceans (largely crayfish and crabs).[43][59]
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To hunt fish, the eagle swoops down over the water and snatches the fish out of the water with its talons. They eat by holding the fish in one claw and tearing the flesh with the other. Eagles have structures on their toes called spicules that allow them to grasp fish. Osprey also have this adaptation.[51] Bald eagles have powerful talons and have been recorded flying with a 6.8 kg (15 lb) mule deer (Odocoileus hemionus) fawn.[99] This feat is the record for the heaviest load carrying ever verified for a flying bird.[100] It has been estimated that the gripping power (pounds by square inch) of the bald eagle is ten times greater than that of a human.[101] Bald eagles can fly with fish at least equal to their own weight, but if the fish is too heavy to lift, the eagle may be dragged into the water. It may swim to safety, in some cases pulling the catch along to the shore as it swims,[102] but some eagles drown or succumb to hypothermia. Many sources claim that bald eagles, like all large eagles, cannot normally take flight carrying prey more than half of their own weight unless aided by favorable wind conditions.[43][71] On numerous occasions, when large prey such as mature salmon or geese are attacked, eagles have been seen to make contact and then drag the prey in a strenuously labored, low flight over the water to a bank, where they then finish off and dismember the prey.[21] When food is abundant, an eagle can gorge itself by storing up to 1 kg (2.2 lb) of food in a pouch in the throat called a crop. Gorging allows the bird to fast for several days if food becomes unavailable.[43] Occasionally, bald eagles may hunt cooperatively when confronting prey, especially relatively large prey such as jackrabbits or herons, with one bird distracting potential prey, while the other comes behind it in order to ambush it.[5][103][104] While hunting waterfowl, bald eagles repeatedly fly at a target and cause it to dive repeatedly, hoping to exhaust the victim so it can be caught (white-tailed eagles have been recorded hunting waterfowl in the same way). When hunting concentrated prey, a successful catch often results in the hunting eagle being pursued by other eagles and needing to find an isolated perch for consumption if it is able to carry it away successfully.[21]
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Unlike some other eagle species, bald eagles rarely take on evasive or dangerous prey on their own. The species mainly target prey which is much smaller than themselves, with most live fish caught weighing 1 to 3 kg (2.2 to 6.6 lb) and most waterbirds preyed weighing 0.2 to 2.7 kg (0.44 to 5.95 lb).[58][71][105] On the other hand, some salmon, carp and marine fish, mammals such as deer fawns and lambs and birds such as swans taken by bald eagles are likely to have been up to at least twice the bald eagles' own size (even if the eagle was unable to fly with it).[19][55][56] They obtain much of their food as carrion or via a practice known as kleptoparasitism, by which they steal prey away from other predators. Due to their dietary habits, bald eagles are frequently viewed in a negative light by humans.[8] Thanks to their superior foraging ability and experience, adults are generally more likely to hunt live prey than immature eagles, which often obtain their food from scavenging.[106][107] They are not very selective about the condition or origin, whether provided by humans, other animals, auto accidents or natural causes, of a carcass's presence, but will avoid eating carrion where disturbances from humans are a regular occurrence. They will scavenge carcasses up to the size of whales, though carcasses of ungulates and large fish are seemingly preferred.[21] Bald eagles also may sometimes feed on material scavenged or stolen from campsites and picnics, as well as garbage dumps (dump usage is habitual mainly in Alaska).[108]
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When competing for food, eagles will usually dominate other fish-eaters and scavengers, aggressively displacing mammals such as coyotes (Canis latrans) and foxes, and birds such as corvids, gulls, vultures and other raptors.[108] Occasionally, coyotes, bobcats (Lynx rufus) and domestic dogs (Canis lupus familiaris) can displace eagles from carrion, usually less confident immature birds, as has been recorded in Maine.[109] Bald eagles are less active, bold predators than golden eagles and get relatively more of their food as carrion and from kleptoparasitism (although it is now generally thought that golden eagles eat more carrion than was previously assumed).[9] However, the two species are roughly equal in size, aggressiveness and physical strength and so competitions can go either way. Neither species is known to be dominant, and the outcome depends on the size and disposition of the individual eagles involved.[21] Wintering bald and golden eagles in Utah both sometimes won conflicts, though in one recorded instance a single bald eagle successfully displaced two consecutive golden eagles from a kill.[110] The bald eagle is thought to be much more numerous in North America than the golden eagle, with the bald species estimated to number at least 150,000 individuals, about twice as many golden eagles there are estimated to live in North America.[9][36] Due to this, bald eagles often outnumber golden eagles at attractive food sources.[9] Despite the potential for contention between these animals, in New Jersey during winter, a golden eagle and numerous bald eagles were observed to hunt snow geese alongside each other without conflict. Similarly, both eagle species have been recorded, via video-monitoring, to feed on gut piles and carcasses of white-tailed deer (Odocoileus virginianus) in remote forest clearings in the eastern Appalachian Mountains without apparent conflict.[9] Bald eagles as frequently mobbed by smaller raptors, due to their infrequent but unpredictable tendency to hunt other birds of prey.[110] Many bald eagles are habitual kleptoparasites, especially in winters when fish are harder to come by. They have been recorded stealing fish from other predators such as ospreys, herons and even otters.[21][111] They have also been recorded opportunistically pirating birds from peregrine falcons (Falco peregrinus), prairie dogs from ferruginous hawks (Buteo regalis) and even jackrabbits from golden eagles.[112][113] When they approach scavengers like dogs, gulls or vultures at carrion sites, they often aggressively attack them and try to force them to disgorge their food.[43] Healthy adult bald eagles are not preyed on in the wild and are thus considered apex predators.[114]
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Bald eagles are sexually mature at four or five years of age. When they are old enough to breed, they often return to the area where they were born. It is thought that bald eagles mate for life. However, if one member of a pair dies or disappears, the survivor will choose a new mate. A pair which has repeatedly failed in breeding attempts may split and look for new mates.[115] Bald eagle courtship involves elaborate, spectacular calls and flight displays. The flight includes swoops, chases, and cartwheels, in which they fly high, lock talons, and free-fall, separating just before hitting the ground.[55][116] Usually, a territory defended by a mature pair will be 1 to 2 km (0.62 to 1.24 mi) of waterside habitat.[8]
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Compared to most other raptors which mostly nest in April or May, bald eagles are early breeders: nest building or reinforcing is often by mid-February, egg laying is often late February (sometimes during deep snow in the North), and incubation is usually mid-March and early May. Eggs hatch from mid April to early May, and the young fledge late June to early July.[8] The nest is the largest of any bird in North America; it is used repeatedly over many years and with new material added each year may eventually be as large as 4 m (13 ft) deep, 2.5 m (8.2 ft) across and weigh 1 metric ton (1.1 short tons);[3] one nest in Florida was found to be 6.1 m (20 ft) deep, 2.9 meters (9.5 ft) across, and to weigh 3 short tons (2.7 metric tons).[117] This nest is on record as the largest tree nest ever recorded for any animal.[118] Usually nests are used for under five years or so, as they either collapse in storms or break the branches supporting them by their sheer weight. However, one nest in the Midwest was occupied continuously for at least 34 years.[43] The nest is built out of branches, usually in large trees found near water. When breeding where there are no trees, the bald eagle will nest on the ground, as has been recorded largely in areas largely isolated from terrestrial predators, such as Amchitka Island in Alaska.[108]
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In Sonora, Mexico, eagles have been observed nesting on top of Hecho catcuses (Pachycereus pectinaboriginum).[119] Nests located on cliffs and rock pinnacles have been reported historically in California, Kansas, Nevada, New Mexico and Utah, but currently are only verified to occur only in Alaska and Arizona.[8] The eggs average about 73 mm (2.9 in) long, ranging from 58 to 85 mm (2.3 to 3.3 in), and have a breadth of 54 mm (2.1 in), ranging from 47 to 63 mm (1.9 to 2.5 in).[51][55] Eggs in Alaska averaged 130 g (4.6 oz) in mass, while in Saskatchewan they averaged 114.4 g (4.04 oz).[120][121] As with their ultimate body size, egg size tends to increase further away from the Equator.[55] Eagles produce between one and three eggs per year, two being typical. Rarely, four eggs have been found in nests but these may be exceptional cases of polygyny.[122] Eagles in captivity have been capable of producing up to seven eggs.[123] It is rare for all three chicks to successfully reach the fledgling stage. The oldest chick often bears the advantage of larger size and louder voice, which tends to draw the parents attention towards it.[8] Occasionally, as is recorded in many large raptorial birds, the oldest sibling sometimes attacks and kills its younger sibling(s), especially early in the nesting period when their sizes are most different.[8] However, nearly half of known bald eagles produce two fledglings (more rarely three), unlike in some other "eagle" species such as some in the genus Aquila, in which a second fledgling is typically observed in less than 20% of nests, despite two eggs typically being laid.[18] Both the male and female take turns incubating the eggs, but the female does most of the sitting. The parent not incubating will hunt for food or look for nesting material during this stage. For the first two to three weeks of the nestling period, at least one adult is at the nest almost 100% of the time. After five to six weeks, the attendance of parents usually drops off considerably (with the parents often perching in trees nearby).[8]
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A young eaglet can gain up to 170 g (6.0 oz) a day, the fastest growth rate of any North American bird.[43] The young eaglets pick up and manipulate sticks, play tug of war with each other, practice holding things in their talons, and stretch and flap their wings. By eight weeks, the eaglets are strong enough to flap their wings, lift their feet off the nest platform, and rise up in the air.[43] The young fledge at anywhere from 8 to 14 weeks of age, though will remain close to the nest and attended to by their parents for a further 6 weeks. Juvenile eagles first start dispersing away from their parents about 8 weeks after they fledge. Variability in departure date related to effects of sex and hatching order on growth and development.[121] For the next four years, immature eagles wander widely in search of food until they attain adult plumage and are eligible to reproduce.[124] Additionally, as shown by a pair of eagles in Shoal Harbor Migratory Bird Sanctuary located near Sidney, British Columbia on June 9, 2017, bald eagles have been recently recorded to occasionally adopt other raptor fledglings into their nests. The pair of eagles in question were recorded carrying a juvenile red-tailed hawk back to their nest, whereupon the chick was accepted into the family by both the parents and the eagles' three fledgelings. Whether or not the chick survived remained to be seen at the time, as young bald eagles are known for killing their siblings. However, the aggression of the red-tailed hawk may ensure its survival, as the hawks are well known for their ability to successfully defend against an eagle attack.[125] Six weeks after however, it was discovered that the hawk, nicknamed "Spunky" by biologists monitoring the nest, had grown to fledgeling size and was learning how to hunt, indicating that it successfully survived.[126]
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The average lifespan of bald eagles in the wild is around 20 years, with the oldest confirmed one having been 38 years of age.[127] In captivity, they often live somewhat longer. In one instance, a captive individual in New York lived for nearly 50 years. As with size, the average lifespan of an eagle population appears to be influenced by its location and access to prey.[128] As they are no longer heavily persecuted, adult mortality is quite low. In one study of Florida eagles, adult bald eagles reportedly had 100% annual survival rate.[9] In Prince William Sound in Alaska, adults had an annual survival rate of 88% even after the Exxon Valdez oil spill adversely affected eagles in the area.[129] Of 1,428 individuals from across the range necropsied by National Wildlife Health Center from 1963 to 1984, 329 (23%) eagles died from trauma, primarily impact with wires and vehicles; 309 (22%) died from gunshot; 158 (11%) died from poisoning; 130 (9%) died from electrocution; 68 (5%) died from trapping; 110 (8%) from emaciation; and 31 (2%) from disease; cause of death was undetermined in 293 (20%) of cases.[130] In this study, 68% of mortality was human-caused.[130] Today eagle-shooting is believed to be considerably reduced due to the species protected status.[131] In one case, an adult eagle investigating a peregrine falcon nest for prey items sustained a concussion from a swooping parent peregrine, and ultimately died days later from it.[132] An early natural history video depicting a cougar (Puma concolor) ambushing and killing an immature bald eagle feeding at a rabbit carcass is viewable online although this film may have been staged.[133]
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Most non-human-related mortality involves nestlings or eggs. Around 50% of eagles survive their first year.[124] However, in the Chesapeake Bay area, 100% of 39 radio-tagged nestlings survived to their first year.[134] Occasionally, nestling or egg fatalities are due to nest collapses, starvation, sibling aggression or inclement weather. Another significant cause of egg and nestling mortality is predation. These have been verified to be preyed by large gulls, corvids (including ravens, crows and magpies), wolverines (Gulo gulo), fishers (Martes pennanti), red-tailed hawks, owls, eagles, bobcats (Lynx rufus), American black bears (Ursus americanus) and raccoons.[120][135][136][137][138][139][140][141] If food access is low, parental attendance at the nest may be lower because both parents may have to forage thus resulting in less protection.[18] Nestlings are usually exempt from predation by terrestrial carnivores that are poor tree-climbers, but Arctic foxes (Vulpes lagopus) occasionally snatched nestlings from ground nests on Amchitka Island in Alaska before they were extirpated from the island.[108] The bald eagle will defend its nest fiercely from all comers and has even repelled attacks from bears, having been recorded knocking a black bear out of a tree when the latter tried to climb a tree holding nestlings.[142]
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Once a common sight in much of the continent, the bald eagle was severely affected in the mid-20th century by a variety of factors, among them the thinning of egg shells attributed to use of the pesticide DDT.[143] Bald eagles, like many birds of prey, were especially affected by DDT due to biomagnification. DDT itself was not lethal to the adult bird, but it interfered with the bird's calcium metabolism, making the bird either sterile or unable to lay healthy eggs. Female eagles laid eggs that were too brittle to withstand the weight of a brooding adult, making it nearly impossible for the eggs to hatch.[34] It is estimated that in the early 18th century, the bald eagle population was 300,000–500,000,[144] but by the 1950s there were only 412 nesting pairs in the 48 contiguous states of the US.[145][146]
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Other factors in bald eagle population reductions were a widespread loss of suitable habitat, as well as both legal and illegal shooting. In 1930 a New York City ornithologist wrote that in the state of Alaska in the previous 12 years approximately 70,000 bald eagles had been shot. Many of the hunters killed the bald eagles under the long-held beliefs that bald eagles grabbed young lambs and even children with their talons, yet the birds were innocent of most of these alleged acts of predation (lamb predation is rare, human predation is thought to be non-existent).[147] Later illegal shooting was described as "the leading cause of direct mortality in both adult and immature bald eagles," according to a 1978 report in the Endangered Species Technical Bulletin. In 1984, the National Wildlife Federation listed hunting, power-line electrocution, and collisions in flight as the leading causes of eagle deaths. Bald eagles have also been killed by oil, lead, and mercury pollution, and by human and predator intrusion at nests.[148]
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The species was first protected in the U.S. and Canada by the 1918 Migratory Bird Treaty, later extended to all of North America. The Bald and Golden Eagle Protection Act, approved by the U.S. Congress in 1940, protected the bald eagle and the golden eagle, prohibiting commercial trapping and killing of the birds. The bald eagle was declared an endangered species in the U.S. in 1967, and amendments to the 1940 act between 1962 and 1972 further restricted commercial uses and increased penalties for violators.[149][150] Perhaps most significant in the species' recovery, in 1972, DDT was banned from usage in the United States due to the fact that it inhibited the reproduction of many birds.[151] DDT was completely banned in Canada in 1989, though its use had been highly restricted since the late 1970s.[152]
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With regulations in place and DDT banned, the eagle population rebounded. The bald eagle can be found in growing concentrations throughout the United States and Canada, particularly near large bodies of water. In the early 1980s, the estimated total population was 100,000 individuals, with 110,000–115,000 by 1992;[3] the U.S. state with the largest resident population is Alaska, with about 40,000–50,000, with the next highest population the Canadian province of British Columbia with 20,000–30,000 in 1992.[3] Obtaining a precise count of bald eagles population is extremely difficult. The most recent data submitted by individual states was in 2006, when 9789 breeding pairs were reported.[153] For some time, the stronghold breeding population of bald eagles in the lower 48 states was in Florida, where over a thousand pairs have held on while populations in other states were significantly reduced by DDT use. Today, the contiguous state with the largest number of breeding pairs of eagles is Minnesota with an estimated 1,312 pairs, surpassing Florida's most recent count of 1,166 pairs. 23, or nearly half, of the 48 contiguous states now have at least 100 breeding pairs of bald eagles.[36] In Washington State, there were only 105 occupied nests in 1980. That number increased by about 30 per year, so that by 2005 there were 840 occupied nests. 2005 was the last year that the Washington Department of Fish and Wildlife counted occupied nests. Further population increases in Washington may be limited by the availability of late winter food, particularly salmon.[154]
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The bald eagle was officially removed from the U.S. federal government's list of endangered species on July 12, 1995, by the U.S. Fish & Wildlife Service, when it was reclassified from "endangered" to "threatened." On July 6, 1999, a proposal was initiated "To Remove the Bald Eagle in the Lower 48 States From the List of Endangered and Threatened Wildlife." It was de-listed on June 28, 2007.[155] It has also been assigned a risk level of least concern category on the IUCN Red List.[2] In the Exxon Valdez oil spill of 1989 an estimated 247 were killed in Prince William Sound, though the local population returned to its pre-spill level by 1995.[5] In some areas, the population has increased such that the eagles are a pest.[156]
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In December 2016, the U.S. Fish and Wildlife Service proposed quadrupling the amount of bald eagles that can be killed by the wind electric generation industry without paying a penalty to 4,200 per year. If issued, the permits would last 30 years, six times the current 5-year permits.[157][158]
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Permits are required to keep bald eagles in captivity in the United States. Permits are primarily issued to public educational institutions, and the eagles which they show are permanently injured individuals which cannot be released to the wild. The facilities where eagles are kept must be equipped with adequate caging and facilities, as well as workers experienced in the handling and care of eagles.[159] The bald eagle can be long-lived in captivity if well cared for, but does not breed well even under the best conditions.[160]
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In Canada[161] and in England[162] a license is required to keep bald eagles for falconry.[163] Bald eagles cannot legally be kept for falconry in the United States, but a license may be issued in some jurisdictions to allow for using such eagles to perform in birds of prey flight shows.[164][165]
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The bald eagle is important in various Native American cultures and, as the national bird of the United States, is prominent in seals and logos, coinage, postage stamps, and other items relating to the U.S. federal government.
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The bald eagle is a sacred bird in some North American cultures, and its feathers, like those of the golden eagle, are central to many religious and spiritual customs among Native Americans. Eagles are considered spiritual messengers between gods and humans by some cultures.[166] Many pow wow dancers use the eagle claw as part of their regalia as well. Eagle feathers are often used in traditional ceremonies, particularly in the construction of regalia worn and as a part of fans, bustles and head dresses. In the Navajo tradition an eagle feather is represented to be a protector, along with the feather Navajo medicine men use the leg and wing bones for ceremonial whistles.[167] The Lakota, for instance, give an eagle feather as a symbol of honor to person who achieves a task. In modern times, it may be given on an event such as a graduation from college.[168] The Pawnee considered eagles as symbols of fertility because their nests are built high off the ground and because they fiercely protect their young.[169] The Choctaw considered the bald eagle, who has direct contact with the upper world of the sun, as a symbol of peace.[170]
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During the Sun Dance, which is practiced by many Plains Indian tribes, the eagle is represented in several ways. The eagle nest is represented by the fork of the lodge where the dance is held. A whistle made from the wing bone of an eagle is used during the course of the dance. Also during the dance, a medicine man may direct his fan, which is made of eagle feathers, to people who seek to be healed. The medicine man touches the fan to the center pole and then to the patient, in order to transmit power from the pole to the patient. The fan is then held up toward the sky, so that the eagle may carry the prayers for the sick to the Creator.[171]
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Current eagle feather law stipulates that only individuals of certifiable Native American ancestry enrolled in a federally recognized tribe are legally authorized to obtain or possess bald or golden eagle feathers for religious or spiritual use. The constitutionality of these laws has been questioned by Native American groups on the basis that it violates the First Amendment by affecting ability to practice their religion freely.[172][173]
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The National Eagle Repository, a division of the FWS, exists as a means to receive, process, and store bald and golden eagles which are found dead, and to distribute the eagles, their parts and feathers, to federally recognized Native American tribes for use in religious ceremonies.[174]
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The bald eagle is the national bird of the United States of America.[175] The founders of the United States were fond of comparing their new republic with the Roman Republic, in which eagle imagery (usually involving the golden eagle) was prominent. On June 20, 1782, the Continental Congress adopted the design for the Great Seal of the United States depicting a bald eagle grasping 13 arrows and an olive branch with thirteen leaves, with its talons.[176][177][178]
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The bald eagle appears on most official seals of the U.S. government, including the presidential seal, the presidential flag, and in the logos of many U.S. federal agencies. Between 1916 and 1945, the presidential flag (but not the seal) showed an eagle facing to its left (the viewer's right), which gave rise to the urban legend that the flag is changed to have the eagle face towards the olive branch in peace, and towards the arrows in wartime.[179]
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Contrary to popular legend, there is no evidence that Benjamin Franklin ever publicly supported the wild turkey (Meleagris gallopavo), rather than the bald eagle, as a symbol of the United States. However, in a letter written to his daughter in 1784 from Paris, criticizing the Society of the Cincinnati, he stated his personal distaste for the bald eagle's behavior. In the letter Franklin states:[180]
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For my own part. I wish the bald eagle had not been chosen the representative of our country. He is a bird of bad moral character. He does not get his living honestly ... besides he is a rank coward: The little king bird not bigger than a sparrow attacks him boldly and drives him out of the district.
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Franklin opposed the creation of the Society because he viewed it, with its hereditary membership, as a noble order unwelcome in the newly independent Republic, contrary to the ideals of Lucius Quinctius Cincinnatus, for whom the Society was named. His reference to the two kinds of birds is interpreted as a satirical comparison between the Society of the Cincinnati and Cincinnatus.[181]
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Largely because of its role as a symbol of the United States, but also because of its being a large predator, the bald eagle has many representations in popular culture. Not all of these representations are accurate. In particular, the movie or television bald eagle typically has a bold, powerful cry. The actual eagle has a much softer, chirpy voice, not in keeping with its popular image; the call of the red-tailed hawk is often substituted in movies and television.[182]
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Pyongyang (US: /ˌpjɒŋˈjæŋ/, UK: /ˌpjʌŋˈjɑːŋ/,[6] Korean: [pʰjʌŋ.jaŋ]) is the capital and largest city of North Korea. Pyongyang is located on the Taedong River about 109 kilometers (68 mi) upstream from its mouth on the Yellow Sea. According to the 2008 population census, it has a population of 3,255,288.[7] Pyongyang is a directly-administered city (직할시; 直轄市; chikhalsi) with equal status to North Korean provinces.
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Pyongyang is considered one of the oldest cities in Korea.[8] It was the capital of two ancient Korean kingdoms, including Gojoseon and Goguryeo, and served as the secondary capital of Goryeo. Much of the city was destroyed during the First Sino-Japanese War, but it was revived under Japanese rule and became an industrial center. Following the establishment of North Korea in 1948, Pyongyang became its de facto capital. The city was again devastated during the Korean War, but was quickly rebuilt after the war with Soviet assistance.
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Pyongyang is the political, industrial and transport center of North Korea. It is home to North Korea's major government institutions, as well as the ruling Workers' Party of Korea.
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The city's other historic names include Kisong, Hwangsong, Rakrang, Sŏgyong, Sodo, Hogyong, Changan,[citation needed] and Heijō[9][10] (during Japanese rule in Korea). There are several variants.[a] During the early 20th century, Pyongyang came to be known among missionaries as being the "Jerusalem of the East", due to its historical status as a stronghold of Christianity, namely Protestantism, especially during the Pyongyang revival of 1907.[22][23]
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After Kim Il-sung's death in 1994, some members of Kim Jong-il's faction proposed changing the name of Pyongyang to "Kim Il-sung City" (Korean: 김일성시; Hanja: 金日成市), but others suggested that North Korea should begin calling Seoul "Kim Il-sung City" instead and grant Pyongyang the moniker "Kim Jong-il City", and in the end neither proposal was implemented.[24]
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In 1955, archaeologists excavated evidence of prehistoric occupation in a large ancient village in the Pyongyang area, called Kŭmtan-ni, dating to the Jeulmun and Mumun pottery periods.[25] North Koreans associate Pyongyang with the mythological city of "Asadal" (Korean: 아사달; Hanja: 阿斯達), or Wanggeom-seong (Korean: 왕검성; Hanja: 王儉城), the first second millennium BC capital of Gojoseon ("Old Joseon") according to Korean historiographies beginning with the 13th-century Samgungnyusa.
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Historians[who?] deny this claim because earlier Chinese historiographical works such as the Guanzi, Classic of Mountains and Seas, Records of the Grand Historian, and Records of the Three Kingdoms, mention a much later "Joseon".[citation needed] The connection between the two therefore may have been asserted by North Korea for the use of propaganda.[citation needed] Nevertheless, Pyongyang became a major city in old Joseon.
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Korean mythology asserts that Pyongyang was founded in 1122 BC on the site of the capital of the legendary king Dangun.[8] Wanggeom-seong, which was in the location of Pyongyang, became the capital of Gojoseon from 194 to 108 BC. It fell in the Han conquest of Gojoseon in 108 BC. Emperor Wu of Han ordered four commanderies be set up, with Lelang Commandery in the center and its capital established as "樂浪" (Old Chinese: *[r]ˤawk*[r]ˤaŋ,[26] Standard Chinese: pinyin: Lèlàng, Korean: 낙랑; Hanja: 樂浪). Several archaeological findings from the later, Eastern Han (20–220 AD) period in the Pyeongyang area seems to suggest that Han forces later launched brief incursions around these parts.
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The area around the city was called Nanglang during the early Three Kingdoms period. As the capital of Nanglang (Korean: 낙랑국; Hanja: 樂浪國),[b] Pyeongyang remained an important commercial and cultural outpost after the Lelang Commandery was destroyed by an expanding Goguryeo in 313.
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Goguryeo moved its capital there in 427. According to Christopher Beckwith, Pyongyang is the Sino-Korean reading of the name they gave it in their language: Piarna, or "level land".[27]
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In 668, Pyongyang became the capital of the Protectorate General to Pacify the East established by the Tang dynasty of China. However, by 676, it was taken by Silla, but left on the border between Silla and Balhae. Pyongyang was left abandoned during the Later Silla period, until it was recovered by Wang Geon and decreed as the Western Capital of Goryeo. During the Joseon period, it became the provincial capital of Pyeongan Province.
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During the Japanese invasions of Korea (1592–98), Pyongyang was captured by the Japanese and held until they were defeated in the Siege of Pyongyang.[8] Later in the 17th century, it became temporarily occupied during the Qing invasion of Joseon until peace arrangements were made between Korea and Qing China. While the invasions made Koreans suspicious of foreigners, the influence of Christianity began to grow after the country opened itself up to foreigners in the 16th century. Pyongyang became the base of Christian expansion in Korea, and by 1880 it had more than 100 churches and more Protestant missionaries than any other Asian city.[8]
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In 1890, the city had 40,000 inhabitants.[28] It was the site of the Battle of Pyongyang during the First Sino-Japanese War, which led to the destruction and depopulation of much of the city. It was the provincial capital of South Pyeongan Province beginning in 1896. Under Japanese colonial rule, the city became an industrial center, called Heijō (with the same Chinese characters 平壤 but read as へいじょう) in Japanese.
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In July 1931, the city experienced anti-Chinese riots as a result of the Wanpaoshan Incident and the sensationalized media reports about it which appeared in Imperial Japanese and Korean newspapers.[29]
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By 1938, Pyongyang had a population of 235,000.[28]
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On 25 August 1945, the Soviet 25th Army entered Pyongyang and it became the temporary capital of the Provisional People's Committee for North Korea. A People's Committee was already established there, led by veteran Christian nationalist Cho Man-sik.[30] Pyongyang became the de facto capital of North Korea upon its establishment in 1948. At the time, the Pyongyang government aimed to recapture Korea's official capital, Seoul. Pyongyang was again severely damaged in the Korean War, during which it was briefly occupied by South Korean forces from 19 October to 6 December 1950. In 1952, it was the target of the largest aerial raid of the entire war, involving 1,400 UN aircraft.
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Already during the war, plans were made to reconstruct the city. On 27 July 1953 – the day the armistice between North Korea and South Korea was signed – The Pyongyang Review wrote: "While streets were in flames, an exhibition showing the general plan of restoration of Pyongyang was held at the Moranbong Underground Theater", the air raid shelter of the government under Moranbong. "On the way of victory... fireworks which streamed high into the night sky of the capital in a gun salute briefly illuminated the construction plan of the city which would rise soon with a new look".[31] After the war, the city was quickly rebuilt with assistance from the Soviet Union, and many buildings were built in the style of Stalinist architecture. The plans for the modern city of Pyongyang were first displayed for public viewing in a theatre building. Kim Jung-hee, one of the founding members of the Korean Architects Alliance, who has studied architecture in prewar Japan was appointed by Kim Il Sung to build the city masterplan, Moscow Architecture Institute, designed the “Pyongyang City Reconstruction and Construction Comprehensive Plan” in 1951, and was officially adopted in 1953. The transformation into a modern, propaganda-designed city called Stalin style architecture with a Korean-style arrangement (and other modernist architecture that was said to have been greatly influenced by Brazilian architect Oscar Niemeyer) began.[32] and in the 1972 Constitution it was officially declared as the capital.
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In 2001, North Korean authorities began a long-term modernisation programme. The Ministry of Capital City Construction Development was included in the Cabinet in that year. In 2006, Kim Jong-il's brother-in-law Jang Song-thaek took charge of the ministry.
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Pyongyang is in the west-central part of North Korea; the city lies on a flat plain about 50 kilometres (31 mi) east of the Korea Bay, an arm of the Yellow Sea. The Taedong River flows southwestward through the city toward the Korea Bay. The Pyongyang plain, where the city is situated, is one of the two large plains on the Western coast of the Korean peninsula, the other being the Chaeryong plain. Both have an area of approximately 500 square kilometers.[33]
|
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Pyongyang has a hot-summer humid continental climate (Köppen: Dwa), featuring warm to hot, humid summers and cold, dry winters.[34][35] Cold, dry winds can blow from Siberia in winter, making conditions very cold; the low temperature is usually below freezing between November and early March, although the average daytime high is at least a few degrees above freezing in every month except January. The winter is generally much drier than summer, with snow falling for 37 days on average.
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The transition from the cold, dry winter to the warm, wet summer occurs rather quickly between April and early May, and there is a similarly abrupt return to winter conditions in late October and November. Summers are generally hot and humid, with the East Asian monsoon taking place from June until September; these are also the hottest months, with average temperatures of 21 to 25 °C (70 to 77 °F), and daytime highs often above 30 °C (86 °F).
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Major government and other public offices are located in Pyongyang, which is constitutionally designated as the country's capital.[39] The seat of the Workers' Party Central Committee and the Pyongyang People's Committee are located in Haebangsan-dong, Chung-guyok. The Cabinet of North Korea is located in Jongro-dong, Chung-guyok.
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Pyongyang is also the seat of all major North Korean security institutions. The largest of them, the Ministry of People's Security, has 130,000 employees working in 12 bureaus. These oversee activities including: police services, security of party officials, classified documents, census, civil registrations, large-scale public construction, traffic control, fire safety, civil defense, public health and customs.[40] Another significant structure based in the city is the State Security Department, whose 30,000 personnel manage intelligence, political prison systems, military industrial security and entry and exit management.[41]
|
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The politics and management of the city is dominated by the Workers' Party of Korea, as they are in the national level. The city is managed by the Pyongyang Party Committee of the Workers' Party of Korea and its chairman is the de facto mayor. The supreme standing state organ is the Pyongyang People's Committee, responsible
|
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+
for everyday events in support of the city. This includes following local Party guidance as channeled through the Pyongyang Party Committee, the distribution of resources prioritised to Pyongyang, and providing support to KWP and internal security agency personnel and families.
|
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P'yŏngyang is divided into 18 wards (ku- or guyŏk) (the city proper) and 2 counties (kun or gun).[42]
|
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|
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+
Foreign media reports in 2010 stated that Kangnam-gun, Chunghwa-gun, Sangwŏn-gun, and Sŭngho-guyŏk had been transferred to the administration of neighboring North Hwanghae province.[43] However, Kangnam-gun was returned to Pyongyang in 2011.[44]
|
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|
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+
After being destroyed during the Korean War, Pyongyang was entirely rebuilt according to Kim Il-sung's vision, which was to create a capital that would boost morale in the post-war years.[45] The result was a city with wide, tree-lined boulevards and public buildings with terraced landscaping, mosaics and decorated ceilings.[46] Its Russian-style architecture makes it reminiscent of a Siberian city during winter snowfall, although edifices of traditional Korean design somewhat soften this perception. In summer, it is notable for its rivers, willow trees, flowers and parkland.[46]
|
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|
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+
The streets are laid out in a north-south, east-west grid, giving the city an orderly appearance.[46] North Korean designers applied the Swedish experience of self-sufficient urban neighbourhoods throughout the entire country, and Pyongyang is no exception. Its inhabitants are mostly divided into administrative units of 5,000 to 6,000 people (dong). These units all have similar sets of amenities including a food store, a barber shop, a tailor, a public bathhouse, a post office, a clinic, a library and others. Many residents occupy high-rise apartment buildings.[47] One of Kim Il-sung's priorities while designing Pyongyang was to limit the population. Authorities maintain a restrictive regime of movement into the city, making it atypical of East Asia as it is silent, uncrowded and spacious.[48]
|
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Structures in Pyongyang are divided into three major architectural categories: monuments, buildings with traditional Korean motifs and high-rises.[49] Some of North Korea's most recognisable landmarks are monuments, like the Juche Tower, the Arch of Triumph and the Mansu Hill Grand Monument. The first of them is a 170-meter granite spire symbolizing the Juche ideology. It was completed in 1982 and contains 25,550 granite blocks, one for each day of Kim Il-sung's life up to that point.[49] The most prominent building on Pyongyang's skyline is Ryugyong Hotel,[49] the seventh highest building in the world terms of floor count, the tallest unoccupied building in the world,[50] and one of the tallest hotels in the world. It has yet to open.[51][52]
|
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Pyongyang has a rapidly evolving skyline, dominated by high-rise apartment buildings. A construction boom began with the Changjon Street Apartment Complex, which was completed in 2012.[53] Construction of the complex began after late leader Kim Jong-il described Changjon Street as "pitiful".[54] Other housing complexes are being upgraded as well, but most are still poorly insulated, and lacking elevators and central heating.[55] An urban renewal program continued under Kim Jong-un's leadership, with the old apartments of the 1970s and '80s replaced by taller high rise buildings and leisure parks like the Kaesong Youth Park, as well as renovations of older buildings.[56] In 2018, the city was described as unrecognizable compared to five years before.[57]
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Notable landmarks in the city include:
|
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Pyongyang TV Tower is a minor landmark. Other visitor attractions include the Korea Central Zoo. The Arch of Reunification has a map of a united Korea supported by two concrete Korean women dressed in traditional dress straddling the Reunification Highway, which stretches from Pyongyang to the Korean Demilitarized Zone (DMZ).
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Juche Tower Monument to the philosophy of Juche (self-reliance)
|
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|
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Arch of Triumph
|
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|
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Arch of Reunification, a monument to the goal of a reunified Korea
|
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Monument to Party Founding
|
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Rungrado May Day Stadium
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Kumsusan Palace of the Sun
|
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Tomb of King Tongmyeong
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Ryugyong Hotel
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Mansu Hill Grand Monument
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Mirae (Future) Scientist's Street
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Pyongyang served as the provincial capital of South Pyongan Province until 1946,[58] and Pyongyang cuisine shares the general culinary tradition of the Pyongan province. The most famous local food is Pyongyang raengmyŏn, or also called mul raengmyŏn or just simply raengmyŏn. Raengmyŏn literally means "cold noodles", while the affix mul refers to water because the dish is served in a cold broth. Raengmyŏn consists of thin and chewy buckwheat noodles in a cold meat-broth with dongchimi (watery kimchi) and topped with a slice of sweet Korean pear.
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Pyongyang raengmyŏn was originally eaten in homes built with ondol (traditional underfloor heating) during the cold winter, so it is also called "Pyongyang deoldeori" (shivering in Pyongyang). Pyongyang locals sometimes enjoyed it as a haejangguk, which is any type of food eaten as a hangover-cure, usually a warm soup.[59]
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Another representative Pyongyang dish, Taedonggang sungeoguk, translates as "trout soup from the Taedong River". The soup features trout (abundant in the Taedong River) along with black peppercorns and salt.[60] Traditionally, it has been served to guests visiting Pyongyang. Therefore, there is a common saying, "How good was the trout soup?", which is used to greet people returning from Pyongyang. Another local specialty, Pyongyang onban (literally "warm rice of Pyongyang") comprises freshly cooked rice topped with sliced mushrooms, chicken, and a couple of bindaetteok (pancakes made from ground mung beans and vegetables).[59]
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In 2018, there were many high quality restaurants in Pyongyang with Korean and international food, and imported alcoholic beverages.[61] Famous restaurants include Okryu-gwan and Ch'ongryugwan.[62] Some street foods exist in Pyongyang, where vendors operate food stalls.[63] Foreign foods like hamburgers, fries, pizza, and coffee are easily found.[61] There is an active nightlife with late-night restaurants and karaoke.[61]
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The city has water parks, amusement parks, skating rinks, health clubs, a shooting range, and a dolphinarium.[56]
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Pyongyang has a number of sports clubs, including the April 25 Sports Club and the Pyongyang City Sports Club.[64] The most popular sport in Pyongyang is football.[citation needed]
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Pyongyang is North Korea's industrial center.[8] Thanks to the abundance of natural resources like coal, iron and limestone, as well as good land and water transport systems, it was the first industrial city to emerge in North Korea after the Korean War. Light and heavy industries are both present and have developed in parallel. Heavy manufactures include cement, industrial ceramics, munitions and weapons, but mechanical engineering remains the core industry. Light industries in Pyongyang and its vicinity include textiles, footwear and food, among others. Special emphasis is put on the production and supply of fresh produce and subsidiary crops in farms on the city's outskirts. Other crops include rice, sweetcorn and soybeans. Pyongyang aims to achieve self-sufficiency in meat production. High-density facilities raise pigs, chicken and other livestock.[8]
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The city still experiences frequent shortages of electricity.[65] To solve this problem, two power stations – Huichon Power Stations 1 and 2 – were built in Chagang Province and supply the city through direct transmission lines. A second phase of the power expansion project was launched in January 2013, consisting of a series of small dams along the Chongchon River. The first two power stations have a maximum generating capacity of 300 megawatts (MW), while the 10 dams to be built under second phase are expected to generate about 120 MW.[65] In addition, the city has several existing or planned thermal power stations. These include Pyongyang TPS with a capacity of 500 MW, East Pyongyang TPS with a capacity of 50 MW, and Kangdong TPS which is under construction.[66]
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Pyongyang is home to several large department stores including the Pothonggang Department Store, Pyongyang Department Store No. 1, Pyongyang Department Store No. 2, Kwangbok Department Store, Ragwon Department Store, Pyongyang Station Department Store, and the Pyongyang Children's Department Store.[67]
|
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The city also has Hwanggumbol Shop, a chain of state-owned convenience stores supplying goods at prices cheaper than those in the jangmadang markets. Hwanggumbol Shops are specifically designed to control North Korea's expanding markets by attracting consumers and guaranteeing the circulation of money in government-operated stores.[68]
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Pyongyang is also the main transport hub of the country: it has a network of roads, railways and air routes which link it to both foreign and domestic destinations. It is the starting point of inter-regional highways reaching Nampo, Wonsan and Kaesong.[8] Pyongyang railway station serves the main railway lines, including the Pyongui Line and the Pyongbu Line. Regular international rail services to Beijing, the Chinese border city of Dandong and Moscow are also available.
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|
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A rail journey to Beijing takes about 25 hours and 25 minutes (K27 from Beijing/K28 from Pyongyang, on Mondays, Wednesdays, Thursdays and Saturdays); a journey to Dandong takes about 6 hours (daily); a journey to Moscow takes six days. The city also connects to the Eurasian Land Bridge via the Trans-Siberian Railway. A high-speed rail link to Wonsan is planned.[69]
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|
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The Metro, tram and trolleybus systems are used mainly by commuters as a primary means of urban transportation.[8] Cycle lanes were introduced on main thoroughfares in July 2015.[70] There are relatively few cars in the city. Cars are a symbol of status in the country due to their scarcity as a result of restrictions on import because of international sanctions and domestic regulations.[71] Some roads are also reported to be in poor condition.[72] However, by 2018, Pyongyang had begun to experience traffic jams.[61]
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State-owned Air Koryo has scheduled international flights from Pyongyang Sunan International Airport to Beijing (PEK), Shenyang (SHE), Vladivostok (VVO), Shanghai (PVG) and Dandong.[73] The only domestic destinations are Hamhung,
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Wonsan, Chongjin, Hyesan and Samjiyon. Since 31 March 2008, Air China launched a regular service between Beijing and Pyongyang,[74] although Air China's flights are often canceled due to lack of passengers.[75]
|
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Kim Il-sung University, North Korea's oldest university, was established in 1946.[8] It has seven colleges, 14 faculties and 16 other institutes, graduate schools and university units.[76] These include the primary medical education and health personnel training unit, the medical college; a physics faculty which covers a range of studies including theoretical physics, optical science, geophysics and astrophysics;[77] an atomic energy institute and a human evolution research office which studies human evolution through a Juche point of view. Kim Il-sung University also has its own publishing house, sports club (Ryongnamsan Sports Club),[78] revolutionary museum, nature museum, libraries, a gym, indoor swimming pool and educator apartment houses. Its two main buildings were completed in 1965 (Building 1) and 1972 (Building 2). A third building on campus is planned.[79]
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Other higher education establishments include Kim Chaek University of Technology, Pyongyang University of Music and Dance and Pyongyang University of Foreign Studies. Pyongyang University of Science and Technology (PUST) is the country's first private university where most of the lecturers are American and courses are carried out in English.[80][81] A science and technology hall is under construction on Ssuk Islet. Its stated purpose is to contribute to the "informatization of educational resources" by centralizing teaching materials, compulsory literature and experimental data for state-level use in a digital format.[82]
|
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Sosong-guyok hosts a 20 MeV cyclotron called MGC-20. The initial project was approved by the International Atomic Energy Agency (IAEA) in 1983 and funded by the IAEA, the United States and the North Korean government. The cyclotron was ordered from the Soviet Union in 1985 and constructed between 1987 and 1990. It is used for student training, production of medical isotopes for nuclear medicine as well as studies in biology, chemistry and physics.[83]
|
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Medical centers include the Red Cross Hospital, the First People's Hospital which is located near Moran Hill and was the first hospital to be built in North Korea after the liberation of Korea in 1945,[84] the Second People's Hospital, Ponghwa Recuperative Center (also known as Bonghwa Clinic or Presidential Clinic) located in Sokam-dong, Potonggang-guyok, 1.5 km (0.93 mi) northwest of Kim Il-sung Square,[85] Pyongyang Medical School Hospital, Namsan Treatment Center which is adjacent[86] Pyongyang's Maternity Hospital, Taesongsan General Hospital,[87] Kim Man-yoo Hospital, Staff Treatment Center and Okryu Children's Hospital. A new hospital named Pyongyang General Hospital is under-construction in Pyongyang.[88]
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Pyongyang is twinned with:
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The Egyptian pyramids are ancient pyramid-shaped masonry structures located in Egypt. As of November 2008, sources cite either 118 or 138 as the number of identified Egyptian pyramids.[1][2] Most were built as tombs for the country's pharaohs and their consorts during the Old and Middle Kingdom periods.[3][4][5]
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The earliest known Egyptian pyramids are found at Saqqara, northwest of Memphis, although at least one step-pyramid-like structure has been found at Saqqara, dating to the 1st Dynasty: Mastaba 3808, which has been attributed to the reign of Pharaoh Anedjib, with inscriptions, and other archaeological remains of the period, suggesting there may have been others.,[6] the otherwise earliest among these is the Pyramid of Djoser built c. 2630–2610 BC during the Third Dynasty.[7] This pyramid and its surrounding complex are generally considered to be the world's oldest monumental structures constructed of dressed masonry.[8]
|
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The most famous Egyptian pyramids are those found at Giza, on the outskirts of Cairo. Several of the Giza pyramids are counted among the largest structures ever built.[9] The Pyramid of Khufu at Giza is the largest Egyptian pyramid. It is the only one of the Seven Wonders of the Ancient World still in existence.[10]
|
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By the time of the Early Dynastic Period, those with sufficient means were buried in bench-like structures known as mastabas.[11][12] At Saqqara, Mastaba 3808, dating from the latter part of the 1st Dynasty, was discovered to contain a large, independently-built step-pyramid-like structure enclosed within the outer palace facade mastaba. Archaeological remains and inscriptions suggest there may have been other similar structures dating to this period.[13]
|
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|
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+
The first historically-documented Egyptian pyramid is attributed by Egyptologists to the 3rd Dynasty pharaoh Djoser. Although Egyptologists often credit his vizier Imhotep as its architect, the dynastic Egyptians themselves, contemporaneously or in numerous later dynastic writings about the character, did not credit him with either designing Djoser's pyramid or the invention of stone architecture.[14] The Pyramid of Djoser was first built as a square mastaba-like structure, which as a rule were known to otherwise be rectangular, and was expanded several times by way of a series of accretion layers, to produce the stepped pyramid structure we see today[15] Egyptologists believe this design served as a gigantic stairway by which the soul of the deceased pharaoh could ascend to the heavens.[16]
|
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Though other pyramids were attempted in the 3rd Dynasty after Djoser, it was the 4th Dynasty, transitioning from the step pyramid to true pyramid shape, which gave rise to the great pyramids of Meidum, Dahsur, and Giza. The last pharaoh of the 4th Dynasty, Shepseskaf, did not build a pyramid and beginning in the 5th Dynasty, for various reasons, the massive scale and precision of construction decreased significantly leaving these later pyramids smaller, less well-built, and often hastily constructed. By the end of the 6th Dynasty pyramid building had largely ended and it was not until the Middle Kingdom that large pyramids were built again, though instead of stone, mud brick was the main construction material.[17]
|
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Long after the end of Egypt's own pyramid-building period, a burst of pyramid building occurred in what is present-day Sudan, after much of Egypt came under the rule of the Kingdom of Kush, which was then based at Napata. Napatan rule, known as the 25th Dynasty, lasted from 750 BCE to 664 BCE, and during that time Egyptian culture made an indelible impression on the Kushites. The Meroitic period of Kushite history, when the kingdom was centered on Meroë, (approximately in the period between 300 BCE and 300 CE), experienced a full-blown pyramid-building revival, which saw more than two hundred Egyptian-inspired indigenous royal pyramid-tombs constructed in the vicinity of the kingdom's capital cities.
|
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Al-Aziz Uthman (1171–1198), the second Ayyubid Sultan of Egypt, tried to destroy the Giza pyramid complex. He gave up after only damaging the Pyramid of Menkaure because the task proved too large.[18]
|
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The shape of Egyptian pyramids is thought to represent the primordial mound from which the Egyptians believed the earth was created. The shape of a pyramid is also thought to be representative of the descending rays of the sun, and most pyramids were faced with polished, highly reflective white limestone, in order to give them a brilliant appearance when viewed from a distance. Pyramids were often also named in ways that referred to solar luminescence. For example, the formal name of the Bent Pyramid at Dahshur was The Southern Shining Pyramid, and that of Senwosret at el-Lahun was Senwosret is Shining.
|
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While it is generally agreed that pyramids were burial monuments, there is continued disagreement on the particular theological principles that might have given rise to them. One suggestion is that they were designed as a type of "resurrection machine."[19]
|
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The Egyptians believed the dark area of the night sky around which the stars appear to revolve was the physical gateway into the heavens. One of the narrow shafts that extend from the main burial chamber through the entire body of the Great Pyramid points directly towards the center of this part of the sky. This suggests the pyramid may have been designed to serve as a means to magically launch the deceased pharaoh's soul directly into the abode of the gods.[19]
|
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All Egyptian pyramids were built on the west bank of the Nile, which, as the site of the setting sun, was associated with the realm of the dead in Egyptian mythology.[20]
|
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+
In 1842, Karl Richard Lepsius produced the first modern list of pyramids – now known as the Lepsius list of pyramids – in which he counted 67. A great many more have since been discovered. As of November 2008, 118 Egyptian pyramids have been identified.[3]
|
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+
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The location of Pyramid 29, which Lepsius called the "Headless Pyramid", was lost for a second time when the structure was buried by desert sands after Lepsius's survey. It was found again only during an archaeological dig conducted in 2008.[21]
|
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Many pyramids are in a poor state of preservation or buried by desert sands. If visible at all, they may appear as little more than mounds of rubble. As a consequence, archaeologists are continuing to identify and study previously unknown pyramid structures.
|
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The most recent pyramid to be discovered was that of Sesheshet at Saqqara, mother of the Sixth Dynasty pharaoh Teti, announced on 11 November 2008.[4][22]
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All of Egypt's pyramids, except the small Third Dynasty pyramid of Zawyet el-Amwat (or Zawyet el-Mayitin), are sited on the west bank of the Nile, and most are grouped together in a number of pyramid fields. The most important of these are listed geographically, from north to south, below.
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Abu Rawash is the site of Egypt's most northerly pyramid (other than the ruins of Lepsius pyramid number one)[5]— the mostly ruined Pyramid of Djedefre, son and successor of Khufu. Originally it was thought that this pyramid had never been completed, but the current archaeological consensus is that not only was it completed, but that it was originally about the same size as the Pyramid of Menkaure, which would have placed it among the half-dozen or so largest pyramids in Egypt.
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Its location adjacent to a major crossroads made it an easy source of stone. Quarrying, which began in Roman times, has left little apart from about 15 courses of stone superimposed upon the natural hillock that formed part of the pyramid's core. A small adjacent satellite pyramid is in a better state of preservation.
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Giza is the location of the Pyramid of Khufu (also known as the "Great Pyramid" and the "Pyramid of Cheops"); the somewhat smaller Pyramid of Khafre (or Chephren); the relatively modest-sized Pyramid of Menkaure (or Mykerinus), along with a number of smaller satellite edifices known as "Queen's pyramids"; and the Great Sphinx of Giza.
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Of the three, only Khafre's pyramid retains part of its original polished limestone casing, near its apex. This pyramid appears larger than the adjacent Khufu pyramid by virtue of its more elevated location, and the steeper angle of inclination of its construction – it is, in fact, smaller in both height and volume.
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The Giza pyramid complex has been a popular tourist destination since antiquity and was popularized in Hellenistic times when the Great Pyramid was listed by Antipater of Sidon as one of the Seven Wonders of the Ancient World. Today it is the only one of those wonders still in existence.
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This site, halfway between Giza and Abusir, is the location for two unfinished Old Kingdom pyramids. The northern structure's owner is believed to be pharaoh Nebka, while the southern structure, known as the Layer Pyramid, may be attributable to the Third Dynasty pharaoh Khaba, a close successor of Sekhemkhet. If this attribution is correct, Khaba's short reign could explain the seemingly unfinished state of this step pyramid. Today it stands around 17 m (56 ft) high; had it been completed, it is likely to have exceeded 40 m (130 ft).
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There are a total of fourteen pyramids at this site, which served as the main royal necropolis during the Fifth Dynasty. The quality of construction of the Abusir pyramids is inferior to those of the Fourth Dynasty – perhaps signaling a decrease in royal power or a less vibrant economy. They are smaller than their predecessors, and are built of low-quality local limestone.
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The three major pyramids are those of Niuserre, which is also the best preserved, Neferirkare Kakai and Sahure. The site is also home to the incomplete Pyramid of Neferefre. Most of the major pyramids at Abusir were built using similar construction techniques, comprising a rubble core surrounded by steps of mud bricks with a limestone outer casing. The largest of these Fifth Dynasty pyramids, the Pyramid of Neferirkare Kakai, is believed to have been built originally as a step pyramid some 70 m (230 ft) high and then later transformed into a "true" pyramid by having its steps filled in with loose masonry.
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Major pyramids located here include the Pyramid of Djoser – generally identified as the world's oldest substantial monumental structure to be built of dressed stone – the Pyramid of Userkaf, the Pyramid of Teti and the Pyramid of Merikare, dating to the First Intermediate Period of Egypt. Also at Saqqara is the Pyramid of Unas, which retains a pyramid causeway that is one of the best-preserved in Egypt. Together with the pyramid of Userkaf, this pyramid was the subject of one of the earliest known restoration attempts, conducted by Khaemweset, a son of Ramesses II.[23] Saqqara is also the location of the incomplete step pyramid of Djoser's successor Sekhemkhet, known as the Buried Pyramid. Archaeologists believe that had this pyramid been completed, it would have been larger than Djoser's.
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South of the main pyramid field at Saqqara is a second collection of later, smaller pyramids, including those of Pepi I, Djedkare Isesi, Merenre, Pepi II and Ibi. Most of these are in a poor state of preservation.
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The Fourth Dynasty pharaoh Shepseskaf either did not share an interest in, or have the capacity to undertake pyramid construction like his predecessors. His tomb, which is also sited at south Saqqara, was instead built as an unusually large mastaba and offering temple complex. It is commonly known as the Mastabat al-Fir’aun.[24]
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A previously unknown pyramid was discovered at north Saqqara in late 2008. Believed to be the tomb of Teti's mother, it currently stands approximately 5 m (16 ft) high, although the original height was closer to 14 m (46 ft).
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This area is arguably the most important pyramid field in Egypt outside Giza and Saqqara, although until 1996 the site was inaccessible due to its location within a military base and was relatively unknown outside archaeological circles.
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The southern Pyramid of Sneferu, commonly known as the Bent Pyramid, is believed to be the first Egyptian pyramid intended by its builders to be a "true" smooth-sided pyramid from the outset; the earlier pyramid at Meidum had smooth sides in its finished state – but it was conceived and built as a step pyramid, before having its steps filled in and concealed beneath a smooth outer casing of dressed stone. As a true smooth-sided structure, the Bent Pyramid was only a partial success – albeit a unique, visually imposing one; it is also the only major Egyptian pyramid to retain a significant proportion of its original smooth outer limestone casing intact. As such it serves as the best contemporary example of how the ancient Egyptians intended their pyramids to look. Several kilometres to the north of the Bent Pyramid is the last – and most successful – of the three pyramids constructed during the reign of Sneferu; the Red Pyramid is the world's first successfully completed smooth-sided pyramid. The structure is also the third largest pyramid in Egypt – after the pyramids of Khufu and Khafra at Giza.
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Also at Dahshur is one of two pyramids built by Amenemhat III, known as the Black Pyramid, as well as a number of small, mostly ruined subsidiary pyramids.
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Located to the south of Dahshur, several mudbrick pyramids were built in this area in the late Middle Kingdom, perhaps for Amenemhat IV and Sobekneferu.
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Two major pyramids are known to have been built at Lisht – those of Amenemhat I and his son, Senusret I. The latter is surrounded by the ruins of ten smaller subsidiary pyramids. One of these subsidiary pyramids is known to be that of Amenemhat's cousin, Khaba II.[25] The site which is in the vicinity of the oasis of the Faiyum, midway between Dahshur and Meidum, and about 100 kilometres south of Cairo, is believed to be in the vicinity of the ancient city of Itjtawy (the precise location of which remains unknown), which served as the capital of Egypt during the Twelfth Dynasty.
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The pyramid at Meidum is one of three constructed during the reign of Sneferu, and is believed by some to have been started by that pharaoh's father and predecessor, Huni. However, that attribution is uncertain, as no record of Huni's name has been found at the site. It was constructed as a step pyramid, and then later converted into the first "true" smooth-sided pyramid when the steps were filled in, and an outer casing added. The pyramid suffered several catastrophic collapses in ancient and medieval times; medieval Arab writers described it as having seven steps – although today only the three uppermost of these remain, giving the structure its odd, tower-like appearance. The hill on which the pyramid is situated is not a natural landscape feature – it is the small mountain of debris created when the lower courses and outer casing of the pyramid gave way.
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Amenemhat III was the last powerful ruler of the Twelfth Dynasty, and the pyramid he built at Hawara, near the Faiyum, is believed to post-date the so-called "Black Pyramid" built by the same ruler at Dahshur. It is the Hawara pyramid that is believed to have been Amenemhet's final resting place.
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The pyramid of Senusret II at el-Lahun is the southernmost royal-tomb pyramid structure in Egypt. Its builders reduced the amount of work necessary to construct it by ingeniously using as its foundation and core a 12-meter-high natural limestone hill.
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Piye, the king of Kush who became the first ruler of the Twenty-fifth Dynasty of Egypt, built a pyramid at El-Kurru. He was the first Egyptian pharaoh to be buried in a pyramid in centuries.
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Taharqa, a Kushite ruler of the Twenty-fifth Dynasty, built his pyramid at Nuri. It was the largest in the area (North Sudan).
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The following table lays out the chronology of the construction of most of the major pyramids mentioned here. Each pyramid is identified through the pharaoh who ordered it built, his approximate reign, and its location.
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*Would have been 91.65 meters (301 feet) or 175 Egyptian Royal cubits.
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*Originally: 143.5 m or 471 feet or 274 Egyptian Royal cubits
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47.6 m (156 ft; 91 Egyptian Royal cubits)
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30 meters (99 feet)
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50 meters (164 feet)
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Constructing the pyramids involved moving huge quantities of stone.
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Papyri discovered at the Egyptian desert near the Red Sea, in 2013 by archaeologist Pierre Tallet, revealed the journal of Merer, an official of Egypt involved in transporting limestone along the Nile River. These papyri reveal processes in the building of the Great Pyramid at Giza, the tomb of the Pharaoh Khufu, just outside modern Cairo.[27]
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Rather than overland transport of the limestone used in building the pyramid, there is evidence that limestone blocks were transported along the Nile River, in the journal of Merer, preserved remnants of ancient canals, and transport boats discovered.[28]
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It is possible that quarried blocks were then transported to the construction site by wooden sleds, with sand in front of the sled wetted to reduce friction. Droplets of water created bridges between the grains of sand, helping them stick together.[29]
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List
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Coordinates: 29°58′34″N 31°07′52″E / 29.97611°N 31.13111°E / 29.97611; 31.13111
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The Great Pyramid of Giza (also known as the Pyramid of Khufu or the Pyramid of Cheops) is the oldest and largest of the three pyramids in the Giza pyramid complex bordering present-day Giza in Greater Cairo, Egypt. It is the oldest of the Seven Wonders of the Ancient World, and the only one to remain largely intact.
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Based on a mark in an interior chamber naming the work gang and a reference to the Fourth Dynasty Egyptian pharaoh Khufu, Egyptologists believe that the pyramid was built as a tomb over a 10- to 20-year period concluding around 2560 BC. Initially standing at 146.5 metres (481 feet), the Great Pyramid was the tallest man-made structure in the world for more than 3,800 years until Lincoln Cathedral was finished in 1311 AD. It is estimated that the pyramid weighs approximately 6 million tonnes, and consists of 2.3 million blocks of limestone and granite, some weighing as much as 80 tonnes. Originally, the Great Pyramid was covered by limestone casing stones that formed a smooth outer surface; what is seen today is the underlying core structure. Some of the casing stones that once covered the structure can still be seen around the base. There have been varying scientific and alternative theories about the Great Pyramid's construction techniques. Most accepted construction hypotheses are based on the idea that it was built by moving huge stones from a quarry and dragging and lifting them into place.
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There are three known chambers inside the Great Pyramid. The lowest chamber is cut into the bedrock upon which the pyramid was built and was unfinished. The so-called[2] Queen's Chamber and King's Chamber are higher up within the pyramid structure. The main part of the Giza complex is a set of buildings that included two mortuary temples in honour of Khufu (one close to the pyramid and one near the Nile), three smaller pyramids for Khufu's wives, an even smaller "satellite" pyramid, a raised causeway connecting the two temples, and small mastaba tombs for nobles surrounding the pyramid.
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Egyptologists believe the pyramid was built as a tomb for the Fourth Dynasty Egyptian pharaoh Khufu (often Hellenized as "Cheops") and was constructed over a 20-year period. Khufu's vizier, Hemiunu (also called Hemon), is believed by some to be the architect of the Great Pyramid.[3] It is thought that, at construction, the Great Pyramid was originally 146.6 metres (481.0 ft) tall, but with the removal of its original casing, its present height is 137 metres (449.5 ft). The lengths of the sides at the base are difficult to reconstruct, given the absence of the casing, but recent analyses put them in a range between 230.26 metres (755.4 ft) and 230.44 metres (756.0 ft). The volume, including an internal hillock, is roughly 2,300,000 cubic metres (81,000,000 cu ft).[4]
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The first precision measurements of the pyramid were made by Egyptologist Sir Flinders Petrie in 1880–82 and published as The Pyramids and Temples of Gizeh.[5] Almost all reports are based on his measurements. Many of the casing-stones and inner chamber blocks of the Great Pyramid fit together with extremely high precision. Based on measurements taken on the north-eastern casing stones, the mean opening of the joints is only 0.5 millimetres (0.020 in) wide.[6]
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The pyramid remained the tallest man-made structure in the world for over 3,800 years,[7] unsurpassed until the 160-metre-tall (520 ft) spire of Lincoln Cathedral was completed c. 1300. The accuracy of the pyramid's workmanship is such that the four sides of the base have an average error of only 58 millimetres in length.[8] The base is horizontal and flat to within ±15 mm (0.6 in).[9] The sides of the square base are closely aligned to the four cardinal compass points (within four minutes of arc)[10] based on true north, not magnetic north,[11] and the finished base was squared to a mean corner error of only 12 seconds of arc.[12]
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The completed design dimensions, as suggested by Petrie's survey and subsequent studies, are estimated to have originally been 280 Egyptian Royal cubits high by 440 cubits long at each of the four sides of its base. The ratio of the perimeter to height of 1760/280 Egyptian Royal cubits equates to 2π to an accuracy of better than 0.05 percent (corresponding to the well-known approximation of π as 22/7). Some Egyptologists consider this to have been the result of deliberate design proportion. Verner wrote, "We can conclude that although the ancient Egyptians could not precisely define the value of π, in practice they used it".[13] Petrie concluded: "but these relations of areas and of circular ratio are so systematic that we should grant that they were in the builder's design".[14] Others have argued that the ancient Egyptians had no concept of pi and would not have thought to encode it in their monuments. They believe that the observed pyramid slope may be based on a simple seked slope choice alone, with no regard to the overall size and proportions of the finished building.[15] In 2013, rolls of papyrus called the Diary of Merer were discovered written by some of those who delivered limestone and other construction materials from Tora to Giza.[16]
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The Great Pyramid consists of an estimated 2.3 million blocks which most believe to have been transported from nearby quarries. The Tura limestone used for the casing was quarried across the river. The largest granite stones in the pyramid, found in the "King's" chamber, weigh 25 to 80 tonnes and were transported from Aswan, more than 800 km (500 mi) away.[citation needed] Ancient Egyptians cut stone into rough blocks by hammering grooves into natural stone faces, inserting wooden wedges, then soaking these with water. As the water was absorbed, the wedges expanded, breaking off workable chunks. Once the blocks were cut, they were carried by boat either up or down the Nile River to the pyramid.[17] It is estimated that 5.5 million tonnes of limestone, 8,000 tonnes of granite (imported from Aswan), and 500,000 tonnes of mortar were used in the construction of the Great Pyramid.[18]
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At completion, the Great Pyramid was surfaced with white "casing stones"—slant-faced, but flat-topped, blocks of highly polished white limestone. These were carefully cut to what is approximately a face slope with a seked of 5+1/2 palms to give the required dimensions. Visibly, all that remains is the underlying stepped core structure seen today.[citation needed] In 1303 AD, a massive earthquake loosened many of the outer casing stones, which in 1356 were carted away by Bahri Sultan An-Nasir Nasir-ad-Din al-Hasan to build mosques and fortresses in nearby Cairo.[citation needed] Many more casing stones were removed from the great pyramids by Muhammad Ali Pasha in the early 19th century to build the upper portion of his Alabaster Mosque in Cairo, not far from Giza.[citation needed] These limestone casings can still be seen as parts of these structures. Later explorers reported massive piles of rubble at the base of the pyramids left over from the continuing collapse of the casing stones, which were subsequently cleared away during continuing excavations of the site.[citation needed]
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Nevertheless, a few of the casing stones from the lowest course can be seen to this day in situ around the base of the Great Pyramid, and display the same workmanship and precision that has been reported for centuries. Petrie also found a different orientation in the core and in the casing measuring 193 centimetres ± 25 centimetres. He suggested a redetermination of north was made after the construction of the core, but a mistake was made, and the casing was built with a different orientation.[5] Petrie related the precision of the casing stones as to being "equal to opticians' work of the present day, but on a scale of acres" and "to place such stones in exact contact would be careful work; but to do so with cement in the joints seems almost impossible".[20] It has been suggested it was the mortar (Petrie's "cement") that made this seemingly impossible task possible, providing a level bed, which enabled the masons to set the stones exactly.[21][22]
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Many alternative, often contradictory, theories have been proposed regarding the pyramid's construction techniques.[23] Many disagree on whether the blocks were dragged, lifted, or even rolled into place. The Greeks believed that slave labour was used, but modern discoveries made at nearby workers' camps associated with construction at Giza suggest that it was built instead by tens of thousands of skilled workers. Verner posited that the labour was organized into a hierarchy, consisting of two gangs of 100,000 men, divided into five zaa or phyle of 20,000 men each, which may have been further divided according to the skills of the workers.[24]
|
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One mystery of the pyramid's construction is its planning. John Romer suggests that they used the same method that had been used for earlier and later constructions, laying out parts of the plan on the ground at a 1-to-1 scale. He writes that "such a working diagram would also serve to generate the architecture of the pyramid with precision unmatched by any other means".[25] He also argues for a 14-year time-span for its construction.[26] A modern construction management study, in association with Mark Lehner and other Egyptologists, estimated that the total project required an average workforce of about 14,500 people and a peak workforce of roughly 40,000. Without the use of pulleys, wheels, or iron tools, they used critical path analysis methods, which suggest that the Great Pyramid was completed from start to finish in approximately 10 years.[27]
|
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The original entrance to the Great Pyramid is on the north, 17 metres (56 ft) vertically above ground level and 7.29 metres (23.9 ft) east of the center line of the pyramid. From this original entrance, there is a Descending Passage 0.96 metres (3.1 ft) high and 1.04 metres (3.4 ft) wide, which goes down at an angle of 26° 31'23" through the masonry of the pyramid and then into the bedrock beneath it. After 105.23 metres (345.2 ft), the passage becomes level and continues for an additional 8.84 metres (29.0 ft) to the lower Chamber, which appears not to have been finished. There is a continuation of the horizontal passage in the south wall of the lower chamber; there is also a pit dug in the floor of the chamber. Some Egyptologists suggest that this Lower Chamber was intended to be the original burial chamber, but Pharaoh Khufu later changed his mind and wanted it to be higher up in the pyramid.[28]
|
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28.2 metres (93 ft) from the entrance is a square hole in the roof of the Descending Passage. Originally concealed with a slab of stone, this is the beginning of the Ascending Passage.[citation needed] The Ascending Passage is 39.3 metres (129 ft) long, as wide and high as the Descending Passage and slopes up at almost precisely the same angle to reach the Grand Gallery. The lower end of the Ascending Passage is closed by three huge blocks of granite, each about 1.5 metres (4.9 ft) long.[citation needed] One must use the Robbers' Tunnel (see below) to access the Ascending Passage.[citation needed] At the start of the Grand Gallery on the right-hand side there is a hole cut in the wall. This is the start of a vertical shaft which follows an irregular path through the masonry of the pyramid to join the Descending Passage. Also at the start of the Grand Gallery there is the Horizontal Passage leading to the "Queen's Chamber". The passage is 1.1m (3'8") high for most of its length, but near the chamber there is a step in the floor, after which the passage is 1.73 metres (5.7 ft) high.[citation needed]
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The "Queen's Chamber"[2] is exactly halfway between the north and south faces of the pyramid and measures 5.75 metres (18.9 ft) north to south, 5.23 metres (17.2 ft) east to west, and has a pointed roof with an apex 6.23 metres (20.4 ft) above the floor. At the eastern end of the chamber there is a niche 4.67 metres (15.3 ft) high. The original depth of the niche was 1.04 metres (3.4 ft), but has since been deepened by treasure hunters.[29]
|
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In the north and south walls of the Queen's Chamber there are shafts, which, unlike those in the King's Chamber that immediately slope upwards (see below), are horizontal for around 2 m (6.6 ft) before sloping upwards. The horizontal distance was cut in 1872 by a British engineer, Waynman Dixon, who believed a shaft similar to those in the King's Chamber must also exist. He was proved right, but because the shafts are not connected to the outer faces of the pyramid or the Queen's Chamber, their purpose is unknown. At the end of one of his shafts, Dixon discovered a ball of black diorite (a type of rock) and a bronze implement of unknown purpose. Both objects are currently in the British Museum.[30]
|
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The shafts in the Queen's Chamber were explored in 1993 by the German engineer Rudolf Gantenbrink using a crawler robot he designed, Upuaut 2. After a climb of 65 m (213 ft),[31] he discovered that one of the shafts was blocked by limestone "doors" with two eroded copper "handles". Some years later the National Geographic Society created a similar robot which, in September 2002, drilled a small hole in the southern door, only to find another door behind it.[32] The northern passage, which was difficult to navigate because of twists and turns, was also found to be blocked by a door.[33]
|
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Research continued in 2011 with the Djedi Project. Realizing the problem was that the National Geographic Society's camera was only able to see straight ahead of it, they instead used a fibre-optic "micro snake camera" that could see around corners. With this they were able to penetrate the first door of the southern shaft through the hole drilled in 2002, and view all the sides of the small chamber behind it. They discovered hieroglyphs written in red paint. They were also able to scrutinize the inside of the two copper "handles" embedded in the door, and they now believe them to be for decorative purposes. They also found the reverse side of the "door" to be finished and polished, which suggests that it was not put there just to block the shaft from debris, but rather for a more specific reason.[34]
|
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|
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The Grand Gallery continues the slope of the Ascending Passage, but is 8.6 metres (28 ft) high and 46.68 metres (153.1 ft) long. At the base it is 2.06 metres (6.8 ft) wide, but after 2.29 metres (7.5 ft) the blocks of stone in the walls are corbelled inwards by 7.6 centimetres (3.0 in) on each side.[citation needed] There are seven of these steps, so, at the top, the Grand Gallery is only 1.04 metres (3.4 ft) wide. It is roofed by slabs of stone laid at a slightly steeper angle than the floor of the gallery, so that each stone fits into a slot cut in the top of the gallery like the teeth of a ratchet. The purpose was to have each block supported by the wall of the Gallery, rather than resting on the block beneath it, in order to prevent cumulative pressure.[35]
|
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|
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At the upper end of the Gallery on the right-hand side there is a hole near the roof that opens into a short tunnel by which access can be gained to the lowest of the Relieving Chambers.[citation needed] The other Relieving Chambers were discovered in 1837–1838 by Colonel Howard Vyse and J.S. Perring, who dug tunnels upwards using blasting powder.[citation needed]
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The floor of the Grand Gallery consists of a shelf or step on either side, 51 centimetres (20 in) wide, leaving a lower ramp 1.04 metres (3.4 ft) wide between them. In the shelves there are 54 slots, 27 on each side matched by vertical and horizontal slots in the walls of the Gallery. These form a cross shape that rises out of the slot in the shelf.[citation needed] The purpose of these slots is not known, but the central gutter in the floor of the Gallery, which is the same width as the Ascending Passage, has led to speculation that the blocking stones were stored in the Grand Gallery and the slots held wooden beams to restrain them from sliding down the passage.[36] This, in turn, has led to the proposal that originally many more than 3 blocking stones were intended, to completely fill the Ascending Passage.[citation needed]
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|
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At the top of the Grand Gallery, there is a step giving onto a horizontal passage some metres long and approximately 1.02 metres (3.3 ft) in height and width, in which can be detected four slots, three of which were probably intended to hold granite portcullises.[citation needed] Fragments of granite found by Petrie in the Descending Passage may have come from these now-vanished doors.[citation needed]
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In 2017, scientists from the ScanPyramids project discovered a large cavity above the Grand Gallery using muon radiography, which they called the "ScanPyramids Big Void". Its length is at least 30 metres (98 ft) and its cross-section is similar to that of the Grand Gallery. Its existence was confirmed by independent detection with three different technologies: nuclear emulsion films, scintillator hodoscopes, and gas detectors.[37][38] The purpose of the cavity is not known and it is not accessible but according to Zahi Hawass it may have been a gap used in the construction of the Grand Gallery.[39] The Japanese research team disputes this, however, saying that the huge void is completely different from the construction spaces previously identified.[40]
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To verify the "ScanPyramids Big Void" and pinpoint the same, a Japanese team of researchers from Kyushu University, Tohoku University, the University of Tokyo and the Chiba Institute of Technology plans to rescan the structure with a newly developed muon detector in 2020.[41]
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The "King's Chamber"[2] is 20 Egyptian Royal cubits or 10.47 metres (34.4 ft) from east to west and 10 cubits or 5.234 metres (17.17 ft) north to south. It has a flat roof 11 cubits and 5 digits or 5.852 metres (19.20 ft) above the floor. 0.91 m (3.0 ft) above the floor there are two narrow shafts in the north and south walls (one is now filled by an extractor fan in an attempt to circulate air inside the pyramid).[citation needed] The purpose of these shafts is not clear: they appear to be aligned towards stars or areas of the northern and southern skies, yet one of them follows a dog-leg course through the masonry, indicating no intention to directly sight stars through them.[citation needed] They were long believed by Egyptologists to be "air shafts" for ventilation, but this idea has now been widely abandoned in favour of the shafts serving a ritualistic purpose associated with the ascension of the king's spirit to the heavens.[42]
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The King's Chamber is entirely faced with granite. Above the roof, which is formed of nine slabs of stone weighing in total about 400 tons, are five compartments known as Relieving Chambers. The first four, like the King's Chamber, have flat roofs formed by the floor of the chamber above, but the final chamber has a pointed roof.[citation needed] Vyse suspected the presence of upper chambers when he found that he could push a long reed through a crack in the ceiling of the first chamber. From lower to upper, the chambers are known as "Davison's Chamber", "Wellington's Chamber", "Nelson's Chamber", "Lady Arbuthnot's Chamber", and "Campbell's Chamber". It is believed that the compartments were intended to safeguard the King's Chamber from the possibility of a roof collapsing under the weight of stone above the Chamber. As the chambers were not intended to be seen, they were not finished in any way and a few of the stones still retain masons' marks painted on them. One of the stones in Campbell's Chamber bears a mark, apparently the name of a work gang.[43][44]
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The only object in the King's Chamber is a rectangular granite sarcophagus, one corner of which is damaged.[citation needed] The sarcophagus is slightly larger than the Ascending Passage, which indicates that it must have been placed in the Chamber before the roof was put in place.[citation needed] Unlike the fine masonry of the walls of the Chamber, the sarcophagus is roughly finished, with saw-marks visible in several places.[citation needed] This is in contrast with the finely finished and decorated sarcophagi found in other pyramids of the same period. Petrie suggested that such a sarcophagus was intended but was lost in the river on the way north from Aswan and a hurriedly made replacement was used instead.[citation needed]
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Today tourists enter the Great Pyramid via the Robbers' Tunnel, which was long ago cut straight through the masonry of the pyramid for approximately 27 metres (89 ft), then turns sharply left to encounter the blocking stones in the Ascending Passage. It is possible to enter the Descending Passage from this point, but access is usually forbidden.[45] The origin of this Robbers' Tunnel is the subject of much scholarly discussion. According to tradition, the chasm was cut around 820 AD by Caliph al-Ma'mun's workmen using a battering ram. According to these accounts, al-Ma'mun's digging dislodged the stone fitted in the ceiling of the Descending Passage to hide the entrance to the Ascending Passage and it was the noise of that stone falling and then sliding down the Descending Passage, which alerted them to the need to turn left. Unable to remove these stones, however, the workmen tunneled up beside them through the softer limestone of the Pyramid until they reached the Ascending Passage.[46][47] Due to a number of historical and archaeological discrepancies, many scholars (with Antoine Isaac Silvestre de Sacy perhaps being the first) contend that this story is apocryphal. They argue that it is much more likely that the tunnel had been carved sometime after the pyramid was initially sealed. This tunnel, the scholars continue, was then resealed (likely during the Ramesside Restoration), and it was this plug that al-Ma'mun's ninth century expedition cleared away.[48]
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The Great Pyramid is surrounded by a complex of several buildings including small pyramids. The Pyramid Temple, which stood on the east side of the pyramid and measured 52.2 metres (171 ft) north to south and 40 metres (130 ft) east to west, has almost entirely disappeared apart from the black basalt paving. There are only a few remnants of the causeway which linked the pyramid with the valley and the Valley Temple. The Valley Temple is buried beneath the village of Nazlet el-Samman; basalt paving and limestone walls have been found but the site has not been excavated.[49][50] The basalt blocks show "clear evidence" of having been cut with some kind of saw with an estimated cutting blade of 15 feet (4.6 m) in length, capable of cutting at a rate of 1.5 inches (38 mm) per minute. Romer suggests that this "super saw" may have had copper teeth and weighed up to 300 pounds (140 kg). He theorizes that such a saw could have been attached to a wooden trestle and possibly used in conjunction with vegetable oil, cutting sand, emery or pounded quartz to cut the blocks, which would have required the labour of at least a dozen men to operate it.[51]
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On the south side are the subsidiary pyramids, popularly known as the Queens' Pyramids. Three remain standing to nearly full height but the fourth was so ruined that its existence was not suspected until the recent discovery of the first course of stones and the remains of the capstone. Hidden beneath the paving around the pyramid was the tomb of Queen Hetepheres I, sister-wife of Sneferu and mother of Khufu. Discovered by accident by the Reisner expedition, the burial was intact, though the carefully sealed coffin proved to be empty.
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A notable construction flanking the Giza pyramid complex is a cyclopean stone wall, the Wall of the Crow.[52] Lehner has discovered a worker's town outside of the wall, otherwise known as "The Lost City", dated by pottery styles, seal impressions, and stratigraphy to have been constructed and occupied sometime during the reigns of Khafre (2520–2494 BC) and Menkaure (2490–2472 BC).[53][54] In the early 21st century, Mark Lehner and his team made several discoveries, including what appears to have been a thriving port, suggesting the town and associated living quarters, which consisted of barracks called "galleries", may not have been for the pyramid workers after all but rather for the soldiers and sailors who utilized the port. In light of this new discovery, as to where then the pyramid workers may have lived, Lehner suggested the alternative possibility they may have camped on the ramps he believes were used to construct the pyramids or possibly at nearby quarries.[55]
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In the early 1970s, the Australian archaeologist Karl Kromer excavated a mound in the South Field of the plateau. This mound contained artefacts including mudbrick seals of Khufu, which he identified with an artisans' settlement.[56] Mudbrick buildings just south of Khufu's Valley Temple contained mud sealings of Khufu and have been suggested to be a settlement serving the cult of Khufu after his death.[57] A worker's cemetery used at least between Khufu's reign and the end of the Fifth Dynasty was discovered south of the Wall of the Crow by Hawass in 1990.[58]
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There are three boat-shaped pits around the pyramid, of a size and shape to have held complete boats, though so shallow that any superstructure, if there ever was one, must have been removed or disassembled. In May 1954, the Egyptian archaeologist Kamal el-Mallakh discovered a fourth pit, a long, narrow rectangle, still covered with slabs of stone weighing up to 15 tons. Inside were 1,224 pieces of wood, the longest 23 metres (75 ft) long, the shortest 10 centimetres (0.33 ft). These were entrusted to a boat builder, Haj Ahmed Yusuf, who worked out how the pieces fit together. The entire process, including conservation and straightening of the warped wood, took fourteen years.
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The result is a cedar-wood boat 43.6 metres (143 ft) long, its timbers held together by ropes, which is currently housed in a special boat-shaped, air-conditioned museum beside the pyramid. During construction of this museum, which stands above the boat pit, a second sealed boat pit was discovered. It was deliberately left unopened until 2011 when excavation began on the boat.[59]
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Although succeeding pyramids were smaller, pyramid-building continued until the end of the Middle Kingdom. However, as authors Brier and Hobbs claim, "all the pyramids were robbed" by the New Kingdom, when the construction of royal tombs in a desert valley, now known as the Valley of the Kings, began.[60][61] Joyce Tyldesley states that the Great Pyramid itself "is known to have been opened and emptied by the Middle Kingdom", before the Arab caliph Al-Ma'mun entered the pyramid around 820 AD.[46]
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I. E. S. Edwards discusses Strabo's mention that the pyramid "a little way up one side has a stone that may be taken out, which being raised up there is a sloping passage to the foundations". Edwards suggested that the pyramid was entered by robbers after the end of the Old Kingdom and sealed and then reopened more than once until Strabo's door was added. He adds: "If this highly speculative surmise be correct, it is also necessary to assume either that the existence of the door was forgotten or that the entrance was again blocked with facing stones", in order to explain why al-Ma'mun could not find the entrance.[62]
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He also discusses a story told by Herodotus. Herodotus visited Egypt in the 5th century BC and recounts a story that he was told concerning vaults under the pyramid built on an island where the body of Cheops lies. Edwards notes that the pyramid had "almost certainly been opened and its contents plundered long before the time of Herodotus" and that it might have been closed again during the Twenty-sixth Dynasty of Egypt when other monuments were restored. He suggests that the story told to Herodotus could have been the result of almost two centuries of telling and retelling by Pyramid guides.[63]
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The Egyptian pyramids are ancient pyramid-shaped masonry structures located in Egypt. As of November 2008, sources cite either 118 or 138 as the number of identified Egyptian pyramids.[1][2] Most were built as tombs for the country's pharaohs and their consorts during the Old and Middle Kingdom periods.[3][4][5]
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The earliest known Egyptian pyramids are found at Saqqara, northwest of Memphis, although at least one step-pyramid-like structure has been found at Saqqara, dating to the 1st Dynasty: Mastaba 3808, which has been attributed to the reign of Pharaoh Anedjib, with inscriptions, and other archaeological remains of the period, suggesting there may have been others.,[6] the otherwise earliest among these is the Pyramid of Djoser built c. 2630–2610 BC during the Third Dynasty.[7] This pyramid and its surrounding complex are generally considered to be the world's oldest monumental structures constructed of dressed masonry.[8]
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The most famous Egyptian pyramids are those found at Giza, on the outskirts of Cairo. Several of the Giza pyramids are counted among the largest structures ever built.[9] The Pyramid of Khufu at Giza is the largest Egyptian pyramid. It is the only one of the Seven Wonders of the Ancient World still in existence.[10]
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By the time of the Early Dynastic Period, those with sufficient means were buried in bench-like structures known as mastabas.[11][12] At Saqqara, Mastaba 3808, dating from the latter part of the 1st Dynasty, was discovered to contain a large, independently-built step-pyramid-like structure enclosed within the outer palace facade mastaba. Archaeological remains and inscriptions suggest there may have been other similar structures dating to this period.[13]
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The first historically-documented Egyptian pyramid is attributed by Egyptologists to the 3rd Dynasty pharaoh Djoser. Although Egyptologists often credit his vizier Imhotep as its architect, the dynastic Egyptians themselves, contemporaneously or in numerous later dynastic writings about the character, did not credit him with either designing Djoser's pyramid or the invention of stone architecture.[14] The Pyramid of Djoser was first built as a square mastaba-like structure, which as a rule were known to otherwise be rectangular, and was expanded several times by way of a series of accretion layers, to produce the stepped pyramid structure we see today[15] Egyptologists believe this design served as a gigantic stairway by which the soul of the deceased pharaoh could ascend to the heavens.[16]
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Though other pyramids were attempted in the 3rd Dynasty after Djoser, it was the 4th Dynasty, transitioning from the step pyramid to true pyramid shape, which gave rise to the great pyramids of Meidum, Dahsur, and Giza. The last pharaoh of the 4th Dynasty, Shepseskaf, did not build a pyramid and beginning in the 5th Dynasty, for various reasons, the massive scale and precision of construction decreased significantly leaving these later pyramids smaller, less well-built, and often hastily constructed. By the end of the 6th Dynasty pyramid building had largely ended and it was not until the Middle Kingdom that large pyramids were built again, though instead of stone, mud brick was the main construction material.[17]
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Long after the end of Egypt's own pyramid-building period, a burst of pyramid building occurred in what is present-day Sudan, after much of Egypt came under the rule of the Kingdom of Kush, which was then based at Napata. Napatan rule, known as the 25th Dynasty, lasted from 750 BCE to 664 BCE, and during that time Egyptian culture made an indelible impression on the Kushites. The Meroitic period of Kushite history, when the kingdom was centered on Meroë, (approximately in the period between 300 BCE and 300 CE), experienced a full-blown pyramid-building revival, which saw more than two hundred Egyptian-inspired indigenous royal pyramid-tombs constructed in the vicinity of the kingdom's capital cities.
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Al-Aziz Uthman (1171–1198), the second Ayyubid Sultan of Egypt, tried to destroy the Giza pyramid complex. He gave up after only damaging the Pyramid of Menkaure because the task proved too large.[18]
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The shape of Egyptian pyramids is thought to represent the primordial mound from which the Egyptians believed the earth was created. The shape of a pyramid is also thought to be representative of the descending rays of the sun, and most pyramids were faced with polished, highly reflective white limestone, in order to give them a brilliant appearance when viewed from a distance. Pyramids were often also named in ways that referred to solar luminescence. For example, the formal name of the Bent Pyramid at Dahshur was The Southern Shining Pyramid, and that of Senwosret at el-Lahun was Senwosret is Shining.
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While it is generally agreed that pyramids were burial monuments, there is continued disagreement on the particular theological principles that might have given rise to them. One suggestion is that they were designed as a type of "resurrection machine."[19]
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The Egyptians believed the dark area of the night sky around which the stars appear to revolve was the physical gateway into the heavens. One of the narrow shafts that extend from the main burial chamber through the entire body of the Great Pyramid points directly towards the center of this part of the sky. This suggests the pyramid may have been designed to serve as a means to magically launch the deceased pharaoh's soul directly into the abode of the gods.[19]
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All Egyptian pyramids were built on the west bank of the Nile, which, as the site of the setting sun, was associated with the realm of the dead in Egyptian mythology.[20]
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In 1842, Karl Richard Lepsius produced the first modern list of pyramids – now known as the Lepsius list of pyramids – in which he counted 67. A great many more have since been discovered. As of November 2008, 118 Egyptian pyramids have been identified.[3]
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The location of Pyramid 29, which Lepsius called the "Headless Pyramid", was lost for a second time when the structure was buried by desert sands after Lepsius's survey. It was found again only during an archaeological dig conducted in 2008.[21]
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Many pyramids are in a poor state of preservation or buried by desert sands. If visible at all, they may appear as little more than mounds of rubble. As a consequence, archaeologists are continuing to identify and study previously unknown pyramid structures.
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The most recent pyramid to be discovered was that of Sesheshet at Saqqara, mother of the Sixth Dynasty pharaoh Teti, announced on 11 November 2008.[4][22]
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All of Egypt's pyramids, except the small Third Dynasty pyramid of Zawyet el-Amwat (or Zawyet el-Mayitin), are sited on the west bank of the Nile, and most are grouped together in a number of pyramid fields. The most important of these are listed geographically, from north to south, below.
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Abu Rawash is the site of Egypt's most northerly pyramid (other than the ruins of Lepsius pyramid number one)[5]— the mostly ruined Pyramid of Djedefre, son and successor of Khufu. Originally it was thought that this pyramid had never been completed, but the current archaeological consensus is that not only was it completed, but that it was originally about the same size as the Pyramid of Menkaure, which would have placed it among the half-dozen or so largest pyramids in Egypt.
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Its location adjacent to a major crossroads made it an easy source of stone. Quarrying, which began in Roman times, has left little apart from about 15 courses of stone superimposed upon the natural hillock that formed part of the pyramid's core. A small adjacent satellite pyramid is in a better state of preservation.
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Giza is the location of the Pyramid of Khufu (also known as the "Great Pyramid" and the "Pyramid of Cheops"); the somewhat smaller Pyramid of Khafre (or Chephren); the relatively modest-sized Pyramid of Menkaure (or Mykerinus), along with a number of smaller satellite edifices known as "Queen's pyramids"; and the Great Sphinx of Giza.
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Of the three, only Khafre's pyramid retains part of its original polished limestone casing, near its apex. This pyramid appears larger than the adjacent Khufu pyramid by virtue of its more elevated location, and the steeper angle of inclination of its construction – it is, in fact, smaller in both height and volume.
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The Giza pyramid complex has been a popular tourist destination since antiquity and was popularized in Hellenistic times when the Great Pyramid was listed by Antipater of Sidon as one of the Seven Wonders of the Ancient World. Today it is the only one of those wonders still in existence.
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This site, halfway between Giza and Abusir, is the location for two unfinished Old Kingdom pyramids. The northern structure's owner is believed to be pharaoh Nebka, while the southern structure, known as the Layer Pyramid, may be attributable to the Third Dynasty pharaoh Khaba, a close successor of Sekhemkhet. If this attribution is correct, Khaba's short reign could explain the seemingly unfinished state of this step pyramid. Today it stands around 17 m (56 ft) high; had it been completed, it is likely to have exceeded 40 m (130 ft).
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There are a total of fourteen pyramids at this site, which served as the main royal necropolis during the Fifth Dynasty. The quality of construction of the Abusir pyramids is inferior to those of the Fourth Dynasty – perhaps signaling a decrease in royal power or a less vibrant economy. They are smaller than their predecessors, and are built of low-quality local limestone.
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The three major pyramids are those of Niuserre, which is also the best preserved, Neferirkare Kakai and Sahure. The site is also home to the incomplete Pyramid of Neferefre. Most of the major pyramids at Abusir were built using similar construction techniques, comprising a rubble core surrounded by steps of mud bricks with a limestone outer casing. The largest of these Fifth Dynasty pyramids, the Pyramid of Neferirkare Kakai, is believed to have been built originally as a step pyramid some 70 m (230 ft) high and then later transformed into a "true" pyramid by having its steps filled in with loose masonry.
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Major pyramids located here include the Pyramid of Djoser – generally identified as the world's oldest substantial monumental structure to be built of dressed stone – the Pyramid of Userkaf, the Pyramid of Teti and the Pyramid of Merikare, dating to the First Intermediate Period of Egypt. Also at Saqqara is the Pyramid of Unas, which retains a pyramid causeway that is one of the best-preserved in Egypt. Together with the pyramid of Userkaf, this pyramid was the subject of one of the earliest known restoration attempts, conducted by Khaemweset, a son of Ramesses II.[23] Saqqara is also the location of the incomplete step pyramid of Djoser's successor Sekhemkhet, known as the Buried Pyramid. Archaeologists believe that had this pyramid been completed, it would have been larger than Djoser's.
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South of the main pyramid field at Saqqara is a second collection of later, smaller pyramids, including those of Pepi I, Djedkare Isesi, Merenre, Pepi II and Ibi. Most of these are in a poor state of preservation.
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The Fourth Dynasty pharaoh Shepseskaf either did not share an interest in, or have the capacity to undertake pyramid construction like his predecessors. His tomb, which is also sited at south Saqqara, was instead built as an unusually large mastaba and offering temple complex. It is commonly known as the Mastabat al-Fir’aun.[24]
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A previously unknown pyramid was discovered at north Saqqara in late 2008. Believed to be the tomb of Teti's mother, it currently stands approximately 5 m (16 ft) high, although the original height was closer to 14 m (46 ft).
|
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|
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This area is arguably the most important pyramid field in Egypt outside Giza and Saqqara, although until 1996 the site was inaccessible due to its location within a military base and was relatively unknown outside archaeological circles.
|
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The southern Pyramid of Sneferu, commonly known as the Bent Pyramid, is believed to be the first Egyptian pyramid intended by its builders to be a "true" smooth-sided pyramid from the outset; the earlier pyramid at Meidum had smooth sides in its finished state – but it was conceived and built as a step pyramid, before having its steps filled in and concealed beneath a smooth outer casing of dressed stone. As a true smooth-sided structure, the Bent Pyramid was only a partial success – albeit a unique, visually imposing one; it is also the only major Egyptian pyramid to retain a significant proportion of its original smooth outer limestone casing intact. As such it serves as the best contemporary example of how the ancient Egyptians intended their pyramids to look. Several kilometres to the north of the Bent Pyramid is the last – and most successful – of the three pyramids constructed during the reign of Sneferu; the Red Pyramid is the world's first successfully completed smooth-sided pyramid. The structure is also the third largest pyramid in Egypt – after the pyramids of Khufu and Khafra at Giza.
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Also at Dahshur is one of two pyramids built by Amenemhat III, known as the Black Pyramid, as well as a number of small, mostly ruined subsidiary pyramids.
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|
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Located to the south of Dahshur, several mudbrick pyramids were built in this area in the late Middle Kingdom, perhaps for Amenemhat IV and Sobekneferu.
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Two major pyramids are known to have been built at Lisht – those of Amenemhat I and his son, Senusret I. The latter is surrounded by the ruins of ten smaller subsidiary pyramids. One of these subsidiary pyramids is known to be that of Amenemhat's cousin, Khaba II.[25] The site which is in the vicinity of the oasis of the Faiyum, midway between Dahshur and Meidum, and about 100 kilometres south of Cairo, is believed to be in the vicinity of the ancient city of Itjtawy (the precise location of which remains unknown), which served as the capital of Egypt during the Twelfth Dynasty.
|
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|
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The pyramid at Meidum is one of three constructed during the reign of Sneferu, and is believed by some to have been started by that pharaoh's father and predecessor, Huni. However, that attribution is uncertain, as no record of Huni's name has been found at the site. It was constructed as a step pyramid, and then later converted into the first "true" smooth-sided pyramid when the steps were filled in, and an outer casing added. The pyramid suffered several catastrophic collapses in ancient and medieval times; medieval Arab writers described it as having seven steps – although today only the three uppermost of these remain, giving the structure its odd, tower-like appearance. The hill on which the pyramid is situated is not a natural landscape feature – it is the small mountain of debris created when the lower courses and outer casing of the pyramid gave way.
|
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Amenemhat III was the last powerful ruler of the Twelfth Dynasty, and the pyramid he built at Hawara, near the Faiyum, is believed to post-date the so-called "Black Pyramid" built by the same ruler at Dahshur. It is the Hawara pyramid that is believed to have been Amenemhet's final resting place.
|
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The pyramid of Senusret II at el-Lahun is the southernmost royal-tomb pyramid structure in Egypt. Its builders reduced the amount of work necessary to construct it by ingeniously using as its foundation and core a 12-meter-high natural limestone hill.
|
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|
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Piye, the king of Kush who became the first ruler of the Twenty-fifth Dynasty of Egypt, built a pyramid at El-Kurru. He was the first Egyptian pharaoh to be buried in a pyramid in centuries.
|
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|
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Taharqa, a Kushite ruler of the Twenty-fifth Dynasty, built his pyramid at Nuri. It was the largest in the area (North Sudan).
|
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|
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The following table lays out the chronology of the construction of most of the major pyramids mentioned here. Each pyramid is identified through the pharaoh who ordered it built, his approximate reign, and its location.
|
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*Would have been 91.65 meters (301 feet) or 175 Egyptian Royal cubits.
|
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*Originally: 143.5 m or 471 feet or 274 Egyptian Royal cubits
|
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|
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47.6 m (156 ft; 91 Egyptian Royal cubits)
|
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30 meters (99 feet)
|
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50 meters (164 feet)
|
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Constructing the pyramids involved moving huge quantities of stone.
|
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Papyri discovered at the Egyptian desert near the Red Sea, in 2013 by archaeologist Pierre Tallet, revealed the journal of Merer, an official of Egypt involved in transporting limestone along the Nile River. These papyri reveal processes in the building of the Great Pyramid at Giza, the tomb of the Pharaoh Khufu, just outside modern Cairo.[27]
|
97 |
+
Rather than overland transport of the limestone used in building the pyramid, there is evidence that limestone blocks were transported along the Nile River, in the journal of Merer, preserved remnants of ancient canals, and transport boats discovered.[28]
|
98 |
+
|
99 |
+
It is possible that quarried blocks were then transported to the construction site by wooden sleds, with sand in front of the sled wetted to reduce friction. Droplets of water created bridges between the grains of sand, helping them stick together.[29]
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List
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Coordinates: 29°58′34″N 31°07′52″E / 29.97611°N 31.13111°E / 29.97611; 31.13111
|
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The Egyptian pyramids are ancient pyramid-shaped masonry structures located in Egypt. As of November 2008, sources cite either 118 or 138 as the number of identified Egyptian pyramids.[1][2] Most were built as tombs for the country's pharaohs and their consorts during the Old and Middle Kingdom periods.[3][4][5]
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The earliest known Egyptian pyramids are found at Saqqara, northwest of Memphis, although at least one step-pyramid-like structure has been found at Saqqara, dating to the 1st Dynasty: Mastaba 3808, which has been attributed to the reign of Pharaoh Anedjib, with inscriptions, and other archaeological remains of the period, suggesting there may have been others.,[6] the otherwise earliest among these is the Pyramid of Djoser built c. 2630–2610 BC during the Third Dynasty.[7] This pyramid and its surrounding complex are generally considered to be the world's oldest monumental structures constructed of dressed masonry.[8]
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|
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The most famous Egyptian pyramids are those found at Giza, on the outskirts of Cairo. Several of the Giza pyramids are counted among the largest structures ever built.[9] The Pyramid of Khufu at Giza is the largest Egyptian pyramid. It is the only one of the Seven Wonders of the Ancient World still in existence.[10]
|
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|
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By the time of the Early Dynastic Period, those with sufficient means were buried in bench-like structures known as mastabas.[11][12] At Saqqara, Mastaba 3808, dating from the latter part of the 1st Dynasty, was discovered to contain a large, independently-built step-pyramid-like structure enclosed within the outer palace facade mastaba. Archaeological remains and inscriptions suggest there may have been other similar structures dating to this period.[13]
|
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|
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The first historically-documented Egyptian pyramid is attributed by Egyptologists to the 3rd Dynasty pharaoh Djoser. Although Egyptologists often credit his vizier Imhotep as its architect, the dynastic Egyptians themselves, contemporaneously or in numerous later dynastic writings about the character, did not credit him with either designing Djoser's pyramid or the invention of stone architecture.[14] The Pyramid of Djoser was first built as a square mastaba-like structure, which as a rule were known to otherwise be rectangular, and was expanded several times by way of a series of accretion layers, to produce the stepped pyramid structure we see today[15] Egyptologists believe this design served as a gigantic stairway by which the soul of the deceased pharaoh could ascend to the heavens.[16]
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|
15 |
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Though other pyramids were attempted in the 3rd Dynasty after Djoser, it was the 4th Dynasty, transitioning from the step pyramid to true pyramid shape, which gave rise to the great pyramids of Meidum, Dahsur, and Giza. The last pharaoh of the 4th Dynasty, Shepseskaf, did not build a pyramid and beginning in the 5th Dynasty, for various reasons, the massive scale and precision of construction decreased significantly leaving these later pyramids smaller, less well-built, and often hastily constructed. By the end of the 6th Dynasty pyramid building had largely ended and it was not until the Middle Kingdom that large pyramids were built again, though instead of stone, mud brick was the main construction material.[17]
|
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|
17 |
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Long after the end of Egypt's own pyramid-building period, a burst of pyramid building occurred in what is present-day Sudan, after much of Egypt came under the rule of the Kingdom of Kush, which was then based at Napata. Napatan rule, known as the 25th Dynasty, lasted from 750 BCE to 664 BCE, and during that time Egyptian culture made an indelible impression on the Kushites. The Meroitic period of Kushite history, when the kingdom was centered on Meroë, (approximately in the period between 300 BCE and 300 CE), experienced a full-blown pyramid-building revival, which saw more than two hundred Egyptian-inspired indigenous royal pyramid-tombs constructed in the vicinity of the kingdom's capital cities.
|
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|
19 |
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Al-Aziz Uthman (1171–1198), the second Ayyubid Sultan of Egypt, tried to destroy the Giza pyramid complex. He gave up after only damaging the Pyramid of Menkaure because the task proved too large.[18]
|
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|
21 |
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The shape of Egyptian pyramids is thought to represent the primordial mound from which the Egyptians believed the earth was created. The shape of a pyramid is also thought to be representative of the descending rays of the sun, and most pyramids were faced with polished, highly reflective white limestone, in order to give them a brilliant appearance when viewed from a distance. Pyramids were often also named in ways that referred to solar luminescence. For example, the formal name of the Bent Pyramid at Dahshur was The Southern Shining Pyramid, and that of Senwosret at el-Lahun was Senwosret is Shining.
|
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|
23 |
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While it is generally agreed that pyramids were burial monuments, there is continued disagreement on the particular theological principles that might have given rise to them. One suggestion is that they were designed as a type of "resurrection machine."[19]
|
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|
25 |
+
The Egyptians believed the dark area of the night sky around which the stars appear to revolve was the physical gateway into the heavens. One of the narrow shafts that extend from the main burial chamber through the entire body of the Great Pyramid points directly towards the center of this part of the sky. This suggests the pyramid may have been designed to serve as a means to magically launch the deceased pharaoh's soul directly into the abode of the gods.[19]
|
26 |
+
|
27 |
+
All Egyptian pyramids were built on the west bank of the Nile, which, as the site of the setting sun, was associated with the realm of the dead in Egyptian mythology.[20]
|
28 |
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|
29 |
+
In 1842, Karl Richard Lepsius produced the first modern list of pyramids – now known as the Lepsius list of pyramids – in which he counted 67. A great many more have since been discovered. As of November 2008, 118 Egyptian pyramids have been identified.[3]
|
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|
31 |
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The location of Pyramid 29, which Lepsius called the "Headless Pyramid", was lost for a second time when the structure was buried by desert sands after Lepsius's survey. It was found again only during an archaeological dig conducted in 2008.[21]
|
32 |
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|
33 |
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Many pyramids are in a poor state of preservation or buried by desert sands. If visible at all, they may appear as little more than mounds of rubble. As a consequence, archaeologists are continuing to identify and study previously unknown pyramid structures.
|
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+
|
35 |
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The most recent pyramid to be discovered was that of Sesheshet at Saqqara, mother of the Sixth Dynasty pharaoh Teti, announced on 11 November 2008.[4][22]
|
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|
37 |
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All of Egypt's pyramids, except the small Third Dynasty pyramid of Zawyet el-Amwat (or Zawyet el-Mayitin), are sited on the west bank of the Nile, and most are grouped together in a number of pyramid fields. The most important of these are listed geographically, from north to south, below.
|
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|
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Abu Rawash is the site of Egypt's most northerly pyramid (other than the ruins of Lepsius pyramid number one)[5]— the mostly ruined Pyramid of Djedefre, son and successor of Khufu. Originally it was thought that this pyramid had never been completed, but the current archaeological consensus is that not only was it completed, but that it was originally about the same size as the Pyramid of Menkaure, which would have placed it among the half-dozen or so largest pyramids in Egypt.
|
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Its location adjacent to a major crossroads made it an easy source of stone. Quarrying, which began in Roman times, has left little apart from about 15 courses of stone superimposed upon the natural hillock that formed part of the pyramid's core. A small adjacent satellite pyramid is in a better state of preservation.
|
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Giza is the location of the Pyramid of Khufu (also known as the "Great Pyramid" and the "Pyramid of Cheops"); the somewhat smaller Pyramid of Khafre (or Chephren); the relatively modest-sized Pyramid of Menkaure (or Mykerinus), along with a number of smaller satellite edifices known as "Queen's pyramids"; and the Great Sphinx of Giza.
|
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Of the three, only Khafre's pyramid retains part of its original polished limestone casing, near its apex. This pyramid appears larger than the adjacent Khufu pyramid by virtue of its more elevated location, and the steeper angle of inclination of its construction – it is, in fact, smaller in both height and volume.
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The Giza pyramid complex has been a popular tourist destination since antiquity and was popularized in Hellenistic times when the Great Pyramid was listed by Antipater of Sidon as one of the Seven Wonders of the Ancient World. Today it is the only one of those wonders still in existence.
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This site, halfway between Giza and Abusir, is the location for two unfinished Old Kingdom pyramids. The northern structure's owner is believed to be pharaoh Nebka, while the southern structure, known as the Layer Pyramid, may be attributable to the Third Dynasty pharaoh Khaba, a close successor of Sekhemkhet. If this attribution is correct, Khaba's short reign could explain the seemingly unfinished state of this step pyramid. Today it stands around 17 m (56 ft) high; had it been completed, it is likely to have exceeded 40 m (130 ft).
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There are a total of fourteen pyramids at this site, which served as the main royal necropolis during the Fifth Dynasty. The quality of construction of the Abusir pyramids is inferior to those of the Fourth Dynasty – perhaps signaling a decrease in royal power or a less vibrant economy. They are smaller than their predecessors, and are built of low-quality local limestone.
|
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+
The three major pyramids are those of Niuserre, which is also the best preserved, Neferirkare Kakai and Sahure. The site is also home to the incomplete Pyramid of Neferefre. Most of the major pyramids at Abusir were built using similar construction techniques, comprising a rubble core surrounded by steps of mud bricks with a limestone outer casing. The largest of these Fifth Dynasty pyramids, the Pyramid of Neferirkare Kakai, is believed to have been built originally as a step pyramid some 70 m (230 ft) high and then later transformed into a "true" pyramid by having its steps filled in with loose masonry.
|
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Major pyramids located here include the Pyramid of Djoser – generally identified as the world's oldest substantial monumental structure to be built of dressed stone – the Pyramid of Userkaf, the Pyramid of Teti and the Pyramid of Merikare, dating to the First Intermediate Period of Egypt. Also at Saqqara is the Pyramid of Unas, which retains a pyramid causeway that is one of the best-preserved in Egypt. Together with the pyramid of Userkaf, this pyramid was the subject of one of the earliest known restoration attempts, conducted by Khaemweset, a son of Ramesses II.[23] Saqqara is also the location of the incomplete step pyramid of Djoser's successor Sekhemkhet, known as the Buried Pyramid. Archaeologists believe that had this pyramid been completed, it would have been larger than Djoser's.
|
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+
|
57 |
+
South of the main pyramid field at Saqqara is a second collection of later, smaller pyramids, including those of Pepi I, Djedkare Isesi, Merenre, Pepi II and Ibi. Most of these are in a poor state of preservation.
|
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+
|
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The Fourth Dynasty pharaoh Shepseskaf either did not share an interest in, or have the capacity to undertake pyramid construction like his predecessors. His tomb, which is also sited at south Saqqara, was instead built as an unusually large mastaba and offering temple complex. It is commonly known as the Mastabat al-Fir’aun.[24]
|
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+
|
61 |
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A previously unknown pyramid was discovered at north Saqqara in late 2008. Believed to be the tomb of Teti's mother, it currently stands approximately 5 m (16 ft) high, although the original height was closer to 14 m (46 ft).
|
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+
|
63 |
+
This area is arguably the most important pyramid field in Egypt outside Giza and Saqqara, although until 1996 the site was inaccessible due to its location within a military base and was relatively unknown outside archaeological circles.
|
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|
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+
The southern Pyramid of Sneferu, commonly known as the Bent Pyramid, is believed to be the first Egyptian pyramid intended by its builders to be a "true" smooth-sided pyramid from the outset; the earlier pyramid at Meidum had smooth sides in its finished state – but it was conceived and built as a step pyramid, before having its steps filled in and concealed beneath a smooth outer casing of dressed stone. As a true smooth-sided structure, the Bent Pyramid was only a partial success – albeit a unique, visually imposing one; it is also the only major Egyptian pyramid to retain a significant proportion of its original smooth outer limestone casing intact. As such it serves as the best contemporary example of how the ancient Egyptians intended their pyramids to look. Several kilometres to the north of the Bent Pyramid is the last – and most successful – of the three pyramids constructed during the reign of Sneferu; the Red Pyramid is the world's first successfully completed smooth-sided pyramid. The structure is also the third largest pyramid in Egypt – after the pyramids of Khufu and Khafra at Giza.
|
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|
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Also at Dahshur is one of two pyramids built by Amenemhat III, known as the Black Pyramid, as well as a number of small, mostly ruined subsidiary pyramids.
|
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|
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Located to the south of Dahshur, several mudbrick pyramids were built in this area in the late Middle Kingdom, perhaps for Amenemhat IV and Sobekneferu.
|
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Two major pyramids are known to have been built at Lisht – those of Amenemhat I and his son, Senusret I. The latter is surrounded by the ruins of ten smaller subsidiary pyramids. One of these subsidiary pyramids is known to be that of Amenemhat's cousin, Khaba II.[25] The site which is in the vicinity of the oasis of the Faiyum, midway between Dahshur and Meidum, and about 100 kilometres south of Cairo, is believed to be in the vicinity of the ancient city of Itjtawy (the precise location of which remains unknown), which served as the capital of Egypt during the Twelfth Dynasty.
|
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|
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The pyramid at Meidum is one of three constructed during the reign of Sneferu, and is believed by some to have been started by that pharaoh's father and predecessor, Huni. However, that attribution is uncertain, as no record of Huni's name has been found at the site. It was constructed as a step pyramid, and then later converted into the first "true" smooth-sided pyramid when the steps were filled in, and an outer casing added. The pyramid suffered several catastrophic collapses in ancient and medieval times; medieval Arab writers described it as having seven steps – although today only the three uppermost of these remain, giving the structure its odd, tower-like appearance. The hill on which the pyramid is situated is not a natural landscape feature – it is the small mountain of debris created when the lower courses and outer casing of the pyramid gave way.
|
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|
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Amenemhat III was the last powerful ruler of the Twelfth Dynasty, and the pyramid he built at Hawara, near the Faiyum, is believed to post-date the so-called "Black Pyramid" built by the same ruler at Dahshur. It is the Hawara pyramid that is believed to have been Amenemhet's final resting place.
|
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|
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The pyramid of Senusret II at el-Lahun is the southernmost royal-tomb pyramid structure in Egypt. Its builders reduced the amount of work necessary to construct it by ingeniously using as its foundation and core a 12-meter-high natural limestone hill.
|
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|
79 |
+
Piye, the king of Kush who became the first ruler of the Twenty-fifth Dynasty of Egypt, built a pyramid at El-Kurru. He was the first Egyptian pharaoh to be buried in a pyramid in centuries.
|
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|
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Taharqa, a Kushite ruler of the Twenty-fifth Dynasty, built his pyramid at Nuri. It was the largest in the area (North Sudan).
|
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|
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The following table lays out the chronology of the construction of most of the major pyramids mentioned here. Each pyramid is identified through the pharaoh who ordered it built, his approximate reign, and its location.
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*Would have been 91.65 meters (301 feet) or 175 Egyptian Royal cubits.
|
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*Originally: 143.5 m or 471 feet or 274 Egyptian Royal cubits
|
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47.6 m (156 ft; 91 Egyptian Royal cubits)
|
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30 meters (99 feet)
|
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50 meters (164 feet)
|
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|
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Constructing the pyramids involved moving huge quantities of stone.
|
96 |
+
Papyri discovered at the Egyptian desert near the Red Sea, in 2013 by archaeologist Pierre Tallet, revealed the journal of Merer, an official of Egypt involved in transporting limestone along the Nile River. These papyri reveal processes in the building of the Great Pyramid at Giza, the tomb of the Pharaoh Khufu, just outside modern Cairo.[27]
|
97 |
+
Rather than overland transport of the limestone used in building the pyramid, there is evidence that limestone blocks were transported along the Nile River, in the journal of Merer, preserved remnants of ancient canals, and transport boats discovered.[28]
|
98 |
+
|
99 |
+
It is possible that quarried blocks were then transported to the construction site by wooden sleds, with sand in front of the sled wetted to reduce friction. Droplets of water created bridges between the grains of sand, helping them stick together.[29]
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List
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Coordinates: 29°58′34″N 31°07′52″E / 29.97611°N 31.13111°E / 29.97611; 31.13111
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The Great Pyramid of Giza (also known as the Pyramid of Khufu or the Pyramid of Cheops) is the oldest and largest of the three pyramids in the Giza pyramid complex bordering present-day Giza in Greater Cairo, Egypt. It is the oldest of the Seven Wonders of the Ancient World, and the only one to remain largely intact.
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Based on a mark in an interior chamber naming the work gang and a reference to the Fourth Dynasty Egyptian pharaoh Khufu, Egyptologists believe that the pyramid was built as a tomb over a 10- to 20-year period concluding around 2560 BC. Initially standing at 146.5 metres (481 feet), the Great Pyramid was the tallest man-made structure in the world for more than 3,800 years until Lincoln Cathedral was finished in 1311 AD. It is estimated that the pyramid weighs approximately 6 million tonnes, and consists of 2.3 million blocks of limestone and granite, some weighing as much as 80 tonnes. Originally, the Great Pyramid was covered by limestone casing stones that formed a smooth outer surface; what is seen today is the underlying core structure. Some of the casing stones that once covered the structure can still be seen around the base. There have been varying scientific and alternative theories about the Great Pyramid's construction techniques. Most accepted construction hypotheses are based on the idea that it was built by moving huge stones from a quarry and dragging and lifting them into place.
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There are three known chambers inside the Great Pyramid. The lowest chamber is cut into the bedrock upon which the pyramid was built and was unfinished. The so-called[2] Queen's Chamber and King's Chamber are higher up within the pyramid structure. The main part of the Giza complex is a set of buildings that included two mortuary temples in honour of Khufu (one close to the pyramid and one near the Nile), three smaller pyramids for Khufu's wives, an even smaller "satellite" pyramid, a raised causeway connecting the two temples, and small mastaba tombs for nobles surrounding the pyramid.
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Egyptologists believe the pyramid was built as a tomb for the Fourth Dynasty Egyptian pharaoh Khufu (often Hellenized as "Cheops") and was constructed over a 20-year period. Khufu's vizier, Hemiunu (also called Hemon), is believed by some to be the architect of the Great Pyramid.[3] It is thought that, at construction, the Great Pyramid was originally 146.6 metres (481.0 ft) tall, but with the removal of its original casing, its present height is 137 metres (449.5 ft). The lengths of the sides at the base are difficult to reconstruct, given the absence of the casing, but recent analyses put them in a range between 230.26 metres (755.4 ft) and 230.44 metres (756.0 ft). The volume, including an internal hillock, is roughly 2,300,000 cubic metres (81,000,000 cu ft).[4]
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The first precision measurements of the pyramid were made by Egyptologist Sir Flinders Petrie in 1880–82 and published as The Pyramids and Temples of Gizeh.[5] Almost all reports are based on his measurements. Many of the casing-stones and inner chamber blocks of the Great Pyramid fit together with extremely high precision. Based on measurements taken on the north-eastern casing stones, the mean opening of the joints is only 0.5 millimetres (0.020 in) wide.[6]
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The pyramid remained the tallest man-made structure in the world for over 3,800 years,[7] unsurpassed until the 160-metre-tall (520 ft) spire of Lincoln Cathedral was completed c. 1300. The accuracy of the pyramid's workmanship is such that the four sides of the base have an average error of only 58 millimetres in length.[8] The base is horizontal and flat to within ±15 mm (0.6 in).[9] The sides of the square base are closely aligned to the four cardinal compass points (within four minutes of arc)[10] based on true north, not magnetic north,[11] and the finished base was squared to a mean corner error of only 12 seconds of arc.[12]
|
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The completed design dimensions, as suggested by Petrie's survey and subsequent studies, are estimated to have originally been 280 Egyptian Royal cubits high by 440 cubits long at each of the four sides of its base. The ratio of the perimeter to height of 1760/280 Egyptian Royal cubits equates to 2π to an accuracy of better than 0.05 percent (corresponding to the well-known approximation of π as 22/7). Some Egyptologists consider this to have been the result of deliberate design proportion. Verner wrote, "We can conclude that although the ancient Egyptians could not precisely define the value of π, in practice they used it".[13] Petrie concluded: "but these relations of areas and of circular ratio are so systematic that we should grant that they were in the builder's design".[14] Others have argued that the ancient Egyptians had no concept of pi and would not have thought to encode it in their monuments. They believe that the observed pyramid slope may be based on a simple seked slope choice alone, with no regard to the overall size and proportions of the finished building.[15] In 2013, rolls of papyrus called the Diary of Merer were discovered written by some of those who delivered limestone and other construction materials from Tora to Giza.[16]
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|
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The Great Pyramid consists of an estimated 2.3 million blocks which most believe to have been transported from nearby quarries. The Tura limestone used for the casing was quarried across the river. The largest granite stones in the pyramid, found in the "King's" chamber, weigh 25 to 80 tonnes and were transported from Aswan, more than 800 km (500 mi) away.[citation needed] Ancient Egyptians cut stone into rough blocks by hammering grooves into natural stone faces, inserting wooden wedges, then soaking these with water. As the water was absorbed, the wedges expanded, breaking off workable chunks. Once the blocks were cut, they were carried by boat either up or down the Nile River to the pyramid.[17] It is estimated that 5.5 million tonnes of limestone, 8,000 tonnes of granite (imported from Aswan), and 500,000 tonnes of mortar were used in the construction of the Great Pyramid.[18]
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At completion, the Great Pyramid was surfaced with white "casing stones"—slant-faced, but flat-topped, blocks of highly polished white limestone. These were carefully cut to what is approximately a face slope with a seked of 5+1/2 palms to give the required dimensions. Visibly, all that remains is the underlying stepped core structure seen today.[citation needed] In 1303 AD, a massive earthquake loosened many of the outer casing stones, which in 1356 were carted away by Bahri Sultan An-Nasir Nasir-ad-Din al-Hasan to build mosques and fortresses in nearby Cairo.[citation needed] Many more casing stones were removed from the great pyramids by Muhammad Ali Pasha in the early 19th century to build the upper portion of his Alabaster Mosque in Cairo, not far from Giza.[citation needed] These limestone casings can still be seen as parts of these structures. Later explorers reported massive piles of rubble at the base of the pyramids left over from the continuing collapse of the casing stones, which were subsequently cleared away during continuing excavations of the site.[citation needed]
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Nevertheless, a few of the casing stones from the lowest course can be seen to this day in situ around the base of the Great Pyramid, and display the same workmanship and precision that has been reported for centuries. Petrie also found a different orientation in the core and in the casing measuring 193 centimetres ± 25 centimetres. He suggested a redetermination of north was made after the construction of the core, but a mistake was made, and the casing was built with a different orientation.[5] Petrie related the precision of the casing stones as to being "equal to opticians' work of the present day, but on a scale of acres" and "to place such stones in exact contact would be careful work; but to do so with cement in the joints seems almost impossible".[20] It has been suggested it was the mortar (Petrie's "cement") that made this seemingly impossible task possible, providing a level bed, which enabled the masons to set the stones exactly.[21][22]
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Many alternative, often contradictory, theories have been proposed regarding the pyramid's construction techniques.[23] Many disagree on whether the blocks were dragged, lifted, or even rolled into place. The Greeks believed that slave labour was used, but modern discoveries made at nearby workers' camps associated with construction at Giza suggest that it was built instead by tens of thousands of skilled workers. Verner posited that the labour was organized into a hierarchy, consisting of two gangs of 100,000 men, divided into five zaa or phyle of 20,000 men each, which may have been further divided according to the skills of the workers.[24]
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One mystery of the pyramid's construction is its planning. John Romer suggests that they used the same method that had been used for earlier and later constructions, laying out parts of the plan on the ground at a 1-to-1 scale. He writes that "such a working diagram would also serve to generate the architecture of the pyramid with precision unmatched by any other means".[25] He also argues for a 14-year time-span for its construction.[26] A modern construction management study, in association with Mark Lehner and other Egyptologists, estimated that the total project required an average workforce of about 14,500 people and a peak workforce of roughly 40,000. Without the use of pulleys, wheels, or iron tools, they used critical path analysis methods, which suggest that the Great Pyramid was completed from start to finish in approximately 10 years.[27]
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The original entrance to the Great Pyramid is on the north, 17 metres (56 ft) vertically above ground level and 7.29 metres (23.9 ft) east of the center line of the pyramid. From this original entrance, there is a Descending Passage 0.96 metres (3.1 ft) high and 1.04 metres (3.4 ft) wide, which goes down at an angle of 26° 31'23" through the masonry of the pyramid and then into the bedrock beneath it. After 105.23 metres (345.2 ft), the passage becomes level and continues for an additional 8.84 metres (29.0 ft) to the lower Chamber, which appears not to have been finished. There is a continuation of the horizontal passage in the south wall of the lower chamber; there is also a pit dug in the floor of the chamber. Some Egyptologists suggest that this Lower Chamber was intended to be the original burial chamber, but Pharaoh Khufu later changed his mind and wanted it to be higher up in the pyramid.[28]
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28.2 metres (93 ft) from the entrance is a square hole in the roof of the Descending Passage. Originally concealed with a slab of stone, this is the beginning of the Ascending Passage.[citation needed] The Ascending Passage is 39.3 metres (129 ft) long, as wide and high as the Descending Passage and slopes up at almost precisely the same angle to reach the Grand Gallery. The lower end of the Ascending Passage is closed by three huge blocks of granite, each about 1.5 metres (4.9 ft) long.[citation needed] One must use the Robbers' Tunnel (see below) to access the Ascending Passage.[citation needed] At the start of the Grand Gallery on the right-hand side there is a hole cut in the wall. This is the start of a vertical shaft which follows an irregular path through the masonry of the pyramid to join the Descending Passage. Also at the start of the Grand Gallery there is the Horizontal Passage leading to the "Queen's Chamber". The passage is 1.1m (3'8") high for most of its length, but near the chamber there is a step in the floor, after which the passage is 1.73 metres (5.7 ft) high.[citation needed]
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The "Queen's Chamber"[2] is exactly halfway between the north and south faces of the pyramid and measures 5.75 metres (18.9 ft) north to south, 5.23 metres (17.2 ft) east to west, and has a pointed roof with an apex 6.23 metres (20.4 ft) above the floor. At the eastern end of the chamber there is a niche 4.67 metres (15.3 ft) high. The original depth of the niche was 1.04 metres (3.4 ft), but has since been deepened by treasure hunters.[29]
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In the north and south walls of the Queen's Chamber there are shafts, which, unlike those in the King's Chamber that immediately slope upwards (see below), are horizontal for around 2 m (6.6 ft) before sloping upwards. The horizontal distance was cut in 1872 by a British engineer, Waynman Dixon, who believed a shaft similar to those in the King's Chamber must also exist. He was proved right, but because the shafts are not connected to the outer faces of the pyramid or the Queen's Chamber, their purpose is unknown. At the end of one of his shafts, Dixon discovered a ball of black diorite (a type of rock) and a bronze implement of unknown purpose. Both objects are currently in the British Museum.[30]
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The shafts in the Queen's Chamber were explored in 1993 by the German engineer Rudolf Gantenbrink using a crawler robot he designed, Upuaut 2. After a climb of 65 m (213 ft),[31] he discovered that one of the shafts was blocked by limestone "doors" with two eroded copper "handles". Some years later the National Geographic Society created a similar robot which, in September 2002, drilled a small hole in the southern door, only to find another door behind it.[32] The northern passage, which was difficult to navigate because of twists and turns, was also found to be blocked by a door.[33]
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Research continued in 2011 with the Djedi Project. Realizing the problem was that the National Geographic Society's camera was only able to see straight ahead of it, they instead used a fibre-optic "micro snake camera" that could see around corners. With this they were able to penetrate the first door of the southern shaft through the hole drilled in 2002, and view all the sides of the small chamber behind it. They discovered hieroglyphs written in red paint. They were also able to scrutinize the inside of the two copper "handles" embedded in the door, and they now believe them to be for decorative purposes. They also found the reverse side of the "door" to be finished and polished, which suggests that it was not put there just to block the shaft from debris, but rather for a more specific reason.[34]
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The Grand Gallery continues the slope of the Ascending Passage, but is 8.6 metres (28 ft) high and 46.68 metres (153.1 ft) long. At the base it is 2.06 metres (6.8 ft) wide, but after 2.29 metres (7.5 ft) the blocks of stone in the walls are corbelled inwards by 7.6 centimetres (3.0 in) on each side.[citation needed] There are seven of these steps, so, at the top, the Grand Gallery is only 1.04 metres (3.4 ft) wide. It is roofed by slabs of stone laid at a slightly steeper angle than the floor of the gallery, so that each stone fits into a slot cut in the top of the gallery like the teeth of a ratchet. The purpose was to have each block supported by the wall of the Gallery, rather than resting on the block beneath it, in order to prevent cumulative pressure.[35]
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At the upper end of the Gallery on the right-hand side there is a hole near the roof that opens into a short tunnel by which access can be gained to the lowest of the Relieving Chambers.[citation needed] The other Relieving Chambers were discovered in 1837–1838 by Colonel Howard Vyse and J.S. Perring, who dug tunnels upwards using blasting powder.[citation needed]
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The floor of the Grand Gallery consists of a shelf or step on either side, 51 centimetres (20 in) wide, leaving a lower ramp 1.04 metres (3.4 ft) wide between them. In the shelves there are 54 slots, 27 on each side matched by vertical and horizontal slots in the walls of the Gallery. These form a cross shape that rises out of the slot in the shelf.[citation needed] The purpose of these slots is not known, but the central gutter in the floor of the Gallery, which is the same width as the Ascending Passage, has led to speculation that the blocking stones were stored in the Grand Gallery and the slots held wooden beams to restrain them from sliding down the passage.[36] This, in turn, has led to the proposal that originally many more than 3 blocking stones were intended, to completely fill the Ascending Passage.[citation needed]
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At the top of the Grand Gallery, there is a step giving onto a horizontal passage some metres long and approximately 1.02 metres (3.3 ft) in height and width, in which can be detected four slots, three of which were probably intended to hold granite portcullises.[citation needed] Fragments of granite found by Petrie in the Descending Passage may have come from these now-vanished doors.[citation needed]
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In 2017, scientists from the ScanPyramids project discovered a large cavity above the Grand Gallery using muon radiography, which they called the "ScanPyramids Big Void". Its length is at least 30 metres (98 ft) and its cross-section is similar to that of the Grand Gallery. Its existence was confirmed by independent detection with three different technologies: nuclear emulsion films, scintillator hodoscopes, and gas detectors.[37][38] The purpose of the cavity is not known and it is not accessible but according to Zahi Hawass it may have been a gap used in the construction of the Grand Gallery.[39] The Japanese research team disputes this, however, saying that the huge void is completely different from the construction spaces previously identified.[40]
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To verify the "ScanPyramids Big Void" and pinpoint the same, a Japanese team of researchers from Kyushu University, Tohoku University, the University of Tokyo and the Chiba Institute of Technology plans to rescan the structure with a newly developed muon detector in 2020.[41]
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The "King's Chamber"[2] is 20 Egyptian Royal cubits or 10.47 metres (34.4 ft) from east to west and 10 cubits or 5.234 metres (17.17 ft) north to south. It has a flat roof 11 cubits and 5 digits or 5.852 metres (19.20 ft) above the floor. 0.91 m (3.0 ft) above the floor there are two narrow shafts in the north and south walls (one is now filled by an extractor fan in an attempt to circulate air inside the pyramid).[citation needed] The purpose of these shafts is not clear: they appear to be aligned towards stars or areas of the northern and southern skies, yet one of them follows a dog-leg course through the masonry, indicating no intention to directly sight stars through them.[citation needed] They were long believed by Egyptologists to be "air shafts" for ventilation, but this idea has now been widely abandoned in favour of the shafts serving a ritualistic purpose associated with the ascension of the king's spirit to the heavens.[42]
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The King's Chamber is entirely faced with granite. Above the roof, which is formed of nine slabs of stone weighing in total about 400 tons, are five compartments known as Relieving Chambers. The first four, like the King's Chamber, have flat roofs formed by the floor of the chamber above, but the final chamber has a pointed roof.[citation needed] Vyse suspected the presence of upper chambers when he found that he could push a long reed through a crack in the ceiling of the first chamber. From lower to upper, the chambers are known as "Davison's Chamber", "Wellington's Chamber", "Nelson's Chamber", "Lady Arbuthnot's Chamber", and "Campbell's Chamber". It is believed that the compartments were intended to safeguard the King's Chamber from the possibility of a roof collapsing under the weight of stone above the Chamber. As the chambers were not intended to be seen, they were not finished in any way and a few of the stones still retain masons' marks painted on them. One of the stones in Campbell's Chamber bears a mark, apparently the name of a work gang.[43][44]
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The only object in the King's Chamber is a rectangular granite sarcophagus, one corner of which is damaged.[citation needed] The sarcophagus is slightly larger than the Ascending Passage, which indicates that it must have been placed in the Chamber before the roof was put in place.[citation needed] Unlike the fine masonry of the walls of the Chamber, the sarcophagus is roughly finished, with saw-marks visible in several places.[citation needed] This is in contrast with the finely finished and decorated sarcophagi found in other pyramids of the same period. Petrie suggested that such a sarcophagus was intended but was lost in the river on the way north from Aswan and a hurriedly made replacement was used instead.[citation needed]
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Today tourists enter the Great Pyramid via the Robbers' Tunnel, which was long ago cut straight through the masonry of the pyramid for approximately 27 metres (89 ft), then turns sharply left to encounter the blocking stones in the Ascending Passage. It is possible to enter the Descending Passage from this point, but access is usually forbidden.[45] The origin of this Robbers' Tunnel is the subject of much scholarly discussion. According to tradition, the chasm was cut around 820 AD by Caliph al-Ma'mun's workmen using a battering ram. According to these accounts, al-Ma'mun's digging dislodged the stone fitted in the ceiling of the Descending Passage to hide the entrance to the Ascending Passage and it was the noise of that stone falling and then sliding down the Descending Passage, which alerted them to the need to turn left. Unable to remove these stones, however, the workmen tunneled up beside them through the softer limestone of the Pyramid until they reached the Ascending Passage.[46][47] Due to a number of historical and archaeological discrepancies, many scholars (with Antoine Isaac Silvestre de Sacy perhaps being the first) contend that this story is apocryphal. They argue that it is much more likely that the tunnel had been carved sometime after the pyramid was initially sealed. This tunnel, the scholars continue, was then resealed (likely during the Ramesside Restoration), and it was this plug that al-Ma'mun's ninth century expedition cleared away.[48]
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The Great Pyramid is surrounded by a complex of several buildings including small pyramids. The Pyramid Temple, which stood on the east side of the pyramid and measured 52.2 metres (171 ft) north to south and 40 metres (130 ft) east to west, has almost entirely disappeared apart from the black basalt paving. There are only a few remnants of the causeway which linked the pyramid with the valley and the Valley Temple. The Valley Temple is buried beneath the village of Nazlet el-Samman; basalt paving and limestone walls have been found but the site has not been excavated.[49][50] The basalt blocks show "clear evidence" of having been cut with some kind of saw with an estimated cutting blade of 15 feet (4.6 m) in length, capable of cutting at a rate of 1.5 inches (38 mm) per minute. Romer suggests that this "super saw" may have had copper teeth and weighed up to 300 pounds (140 kg). He theorizes that such a saw could have been attached to a wooden trestle and possibly used in conjunction with vegetable oil, cutting sand, emery or pounded quartz to cut the blocks, which would have required the labour of at least a dozen men to operate it.[51]
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On the south side are the subsidiary pyramids, popularly known as the Queens' Pyramids. Three remain standing to nearly full height but the fourth was so ruined that its existence was not suspected until the recent discovery of the first course of stones and the remains of the capstone. Hidden beneath the paving around the pyramid was the tomb of Queen Hetepheres I, sister-wife of Sneferu and mother of Khufu. Discovered by accident by the Reisner expedition, the burial was intact, though the carefully sealed coffin proved to be empty.
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A notable construction flanking the Giza pyramid complex is a cyclopean stone wall, the Wall of the Crow.[52] Lehner has discovered a worker's town outside of the wall, otherwise known as "The Lost City", dated by pottery styles, seal impressions, and stratigraphy to have been constructed and occupied sometime during the reigns of Khafre (2520–2494 BC) and Menkaure (2490–2472 BC).[53][54] In the early 21st century, Mark Lehner and his team made several discoveries, including what appears to have been a thriving port, suggesting the town and associated living quarters, which consisted of barracks called "galleries", may not have been for the pyramid workers after all but rather for the soldiers and sailors who utilized the port. In light of this new discovery, as to where then the pyramid workers may have lived, Lehner suggested the alternative possibility they may have camped on the ramps he believes were used to construct the pyramids or possibly at nearby quarries.[55]
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In the early 1970s, the Australian archaeologist Karl Kromer excavated a mound in the South Field of the plateau. This mound contained artefacts including mudbrick seals of Khufu, which he identified with an artisans' settlement.[56] Mudbrick buildings just south of Khufu's Valley Temple contained mud sealings of Khufu and have been suggested to be a settlement serving the cult of Khufu after his death.[57] A worker's cemetery used at least between Khufu's reign and the end of the Fifth Dynasty was discovered south of the Wall of the Crow by Hawass in 1990.[58]
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There are three boat-shaped pits around the pyramid, of a size and shape to have held complete boats, though so shallow that any superstructure, if there ever was one, must have been removed or disassembled. In May 1954, the Egyptian archaeologist Kamal el-Mallakh discovered a fourth pit, a long, narrow rectangle, still covered with slabs of stone weighing up to 15 tons. Inside were 1,224 pieces of wood, the longest 23 metres (75 ft) long, the shortest 10 centimetres (0.33 ft). These were entrusted to a boat builder, Haj Ahmed Yusuf, who worked out how the pieces fit together. The entire process, including conservation and straightening of the warped wood, took fourteen years.
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The result is a cedar-wood boat 43.6 metres (143 ft) long, its timbers held together by ropes, which is currently housed in a special boat-shaped, air-conditioned museum beside the pyramid. During construction of this museum, which stands above the boat pit, a second sealed boat pit was discovered. It was deliberately left unopened until 2011 when excavation began on the boat.[59]
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Although succeeding pyramids were smaller, pyramid-building continued until the end of the Middle Kingdom. However, as authors Brier and Hobbs claim, "all the pyramids were robbed" by the New Kingdom, when the construction of royal tombs in a desert valley, now known as the Valley of the Kings, began.[60][61] Joyce Tyldesley states that the Great Pyramid itself "is known to have been opened and emptied by the Middle Kingdom", before the Arab caliph Al-Ma'mun entered the pyramid around 820 AD.[46]
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I. E. S. Edwards discusses Strabo's mention that the pyramid "a little way up one side has a stone that may be taken out, which being raised up there is a sloping passage to the foundations". Edwards suggested that the pyramid was entered by robbers after the end of the Old Kingdom and sealed and then reopened more than once until Strabo's door was added. He adds: "If this highly speculative surmise be correct, it is also necessary to assume either that the existence of the door was forgotten or that the entrance was again blocked with facing stones", in order to explain why al-Ma'mun could not find the entrance.[62]
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He also discusses a story told by Herodotus. Herodotus visited Egypt in the 5th century BC and recounts a story that he was told concerning vaults under the pyramid built on an island where the body of Cheops lies. Edwards notes that the pyramid had "almost certainly been opened and its contents plundered long before the time of Herodotus" and that it might have been closed again during the Twenty-sixth Dynasty of Egypt when other monuments were restored. He suggests that the story told to Herodotus could have been the result of almost two centuries of telling and retelling by Pyramid guides.[63]
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The Egyptian pyramids are ancient pyramid-shaped masonry structures located in Egypt. As of November 2008, sources cite either 118 or 138 as the number of identified Egyptian pyramids.[1][2] Most were built as tombs for the country's pharaohs and their consorts during the Old and Middle Kingdom periods.[3][4][5]
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The earliest known Egyptian pyramids are found at Saqqara, northwest of Memphis, although at least one step-pyramid-like structure has been found at Saqqara, dating to the 1st Dynasty: Mastaba 3808, which has been attributed to the reign of Pharaoh Anedjib, with inscriptions, and other archaeological remains of the period, suggesting there may have been others.,[6] the otherwise earliest among these is the Pyramid of Djoser built c. 2630–2610 BC during the Third Dynasty.[7] This pyramid and its surrounding complex are generally considered to be the world's oldest monumental structures constructed of dressed masonry.[8]
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The most famous Egyptian pyramids are those found at Giza, on the outskirts of Cairo. Several of the Giza pyramids are counted among the largest structures ever built.[9] The Pyramid of Khufu at Giza is the largest Egyptian pyramid. It is the only one of the Seven Wonders of the Ancient World still in existence.[10]
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By the time of the Early Dynastic Period, those with sufficient means were buried in bench-like structures known as mastabas.[11][12] At Saqqara, Mastaba 3808, dating from the latter part of the 1st Dynasty, was discovered to contain a large, independently-built step-pyramid-like structure enclosed within the outer palace facade mastaba. Archaeological remains and inscriptions suggest there may have been other similar structures dating to this period.[13]
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The first historically-documented Egyptian pyramid is attributed by Egyptologists to the 3rd Dynasty pharaoh Djoser. Although Egyptologists often credit his vizier Imhotep as its architect, the dynastic Egyptians themselves, contemporaneously or in numerous later dynastic writings about the character, did not credit him with either designing Djoser's pyramid or the invention of stone architecture.[14] The Pyramid of Djoser was first built as a square mastaba-like structure, which as a rule were known to otherwise be rectangular, and was expanded several times by way of a series of accretion layers, to produce the stepped pyramid structure we see today[15] Egyptologists believe this design served as a gigantic stairway by which the soul of the deceased pharaoh could ascend to the heavens.[16]
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Though other pyramids were attempted in the 3rd Dynasty after Djoser, it was the 4th Dynasty, transitioning from the step pyramid to true pyramid shape, which gave rise to the great pyramids of Meidum, Dahsur, and Giza. The last pharaoh of the 4th Dynasty, Shepseskaf, did not build a pyramid and beginning in the 5th Dynasty, for various reasons, the massive scale and precision of construction decreased significantly leaving these later pyramids smaller, less well-built, and often hastily constructed. By the end of the 6th Dynasty pyramid building had largely ended and it was not until the Middle Kingdom that large pyramids were built again, though instead of stone, mud brick was the main construction material.[17]
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Long after the end of Egypt's own pyramid-building period, a burst of pyramid building occurred in what is present-day Sudan, after much of Egypt came under the rule of the Kingdom of Kush, which was then based at Napata. Napatan rule, known as the 25th Dynasty, lasted from 750 BCE to 664 BCE, and during that time Egyptian culture made an indelible impression on the Kushites. The Meroitic period of Kushite history, when the kingdom was centered on Meroë, (approximately in the period between 300 BCE and 300 CE), experienced a full-blown pyramid-building revival, which saw more than two hundred Egyptian-inspired indigenous royal pyramid-tombs constructed in the vicinity of the kingdom's capital cities.
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Al-Aziz Uthman (1171–1198), the second Ayyubid Sultan of Egypt, tried to destroy the Giza pyramid complex. He gave up after only damaging the Pyramid of Menkaure because the task proved too large.[18]
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The shape of Egyptian pyramids is thought to represent the primordial mound from which the Egyptians believed the earth was created. The shape of a pyramid is also thought to be representative of the descending rays of the sun, and most pyramids were faced with polished, highly reflective white limestone, in order to give them a brilliant appearance when viewed from a distance. Pyramids were often also named in ways that referred to solar luminescence. For example, the formal name of the Bent Pyramid at Dahshur was The Southern Shining Pyramid, and that of Senwosret at el-Lahun was Senwosret is Shining.
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While it is generally agreed that pyramids were burial monuments, there is continued disagreement on the particular theological principles that might have given rise to them. One suggestion is that they were designed as a type of "resurrection machine."[19]
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The Egyptians believed the dark area of the night sky around which the stars appear to revolve was the physical gateway into the heavens. One of the narrow shafts that extend from the main burial chamber through the entire body of the Great Pyramid points directly towards the center of this part of the sky. This suggests the pyramid may have been designed to serve as a means to magically launch the deceased pharaoh's soul directly into the abode of the gods.[19]
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All Egyptian pyramids were built on the west bank of the Nile, which, as the site of the setting sun, was associated with the realm of the dead in Egyptian mythology.[20]
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In 1842, Karl Richard Lepsius produced the first modern list of pyramids – now known as the Lepsius list of pyramids – in which he counted 67. A great many more have since been discovered. As of November 2008, 118 Egyptian pyramids have been identified.[3]
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The location of Pyramid 29, which Lepsius called the "Headless Pyramid", was lost for a second time when the structure was buried by desert sands after Lepsius's survey. It was found again only during an archaeological dig conducted in 2008.[21]
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Many pyramids are in a poor state of preservation or buried by desert sands. If visible at all, they may appear as little more than mounds of rubble. As a consequence, archaeologists are continuing to identify and study previously unknown pyramid structures.
|
34 |
+
|
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+
The most recent pyramid to be discovered was that of Sesheshet at Saqqara, mother of the Sixth Dynasty pharaoh Teti, announced on 11 November 2008.[4][22]
|
36 |
+
|
37 |
+
All of Egypt's pyramids, except the small Third Dynasty pyramid of Zawyet el-Amwat (or Zawyet el-Mayitin), are sited on the west bank of the Nile, and most are grouped together in a number of pyramid fields. The most important of these are listed geographically, from north to south, below.
|
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|
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Abu Rawash is the site of Egypt's most northerly pyramid (other than the ruins of Lepsius pyramid number one)[5]— the mostly ruined Pyramid of Djedefre, son and successor of Khufu. Originally it was thought that this pyramid had never been completed, but the current archaeological consensus is that not only was it completed, but that it was originally about the same size as the Pyramid of Menkaure, which would have placed it among the half-dozen or so largest pyramids in Egypt.
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Its location adjacent to a major crossroads made it an easy source of stone. Quarrying, which began in Roman times, has left little apart from about 15 courses of stone superimposed upon the natural hillock that formed part of the pyramid's core. A small adjacent satellite pyramid is in a better state of preservation.
|
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|
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Giza is the location of the Pyramid of Khufu (also known as the "Great Pyramid" and the "Pyramid of Cheops"); the somewhat smaller Pyramid of Khafre (or Chephren); the relatively modest-sized Pyramid of Menkaure (or Mykerinus), along with a number of smaller satellite edifices known as "Queen's pyramids"; and the Great Sphinx of Giza.
|
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Of the three, only Khafre's pyramid retains part of its original polished limestone casing, near its apex. This pyramid appears larger than the adjacent Khufu pyramid by virtue of its more elevated location, and the steeper angle of inclination of its construction – it is, in fact, smaller in both height and volume.
|
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The Giza pyramid complex has been a popular tourist destination since antiquity and was popularized in Hellenistic times when the Great Pyramid was listed by Antipater of Sidon as one of the Seven Wonders of the Ancient World. Today it is the only one of those wonders still in existence.
|
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|
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This site, halfway between Giza and Abusir, is the location for two unfinished Old Kingdom pyramids. The northern structure's owner is believed to be pharaoh Nebka, while the southern structure, known as the Layer Pyramid, may be attributable to the Third Dynasty pharaoh Khaba, a close successor of Sekhemkhet. If this attribution is correct, Khaba's short reign could explain the seemingly unfinished state of this step pyramid. Today it stands around 17 m (56 ft) high; had it been completed, it is likely to have exceeded 40 m (130 ft).
|
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|
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There are a total of fourteen pyramids at this site, which served as the main royal necropolis during the Fifth Dynasty. The quality of construction of the Abusir pyramids is inferior to those of the Fourth Dynasty – perhaps signaling a decrease in royal power or a less vibrant economy. They are smaller than their predecessors, and are built of low-quality local limestone.
|
52 |
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|
53 |
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The three major pyramids are those of Niuserre, which is also the best preserved, Neferirkare Kakai and Sahure. The site is also home to the incomplete Pyramid of Neferefre. Most of the major pyramids at Abusir were built using similar construction techniques, comprising a rubble core surrounded by steps of mud bricks with a limestone outer casing. The largest of these Fifth Dynasty pyramids, the Pyramid of Neferirkare Kakai, is believed to have been built originally as a step pyramid some 70 m (230 ft) high and then later transformed into a "true" pyramid by having its steps filled in with loose masonry.
|
54 |
+
|
55 |
+
Major pyramids located here include the Pyramid of Djoser – generally identified as the world's oldest substantial monumental structure to be built of dressed stone – the Pyramid of Userkaf, the Pyramid of Teti and the Pyramid of Merikare, dating to the First Intermediate Period of Egypt. Also at Saqqara is the Pyramid of Unas, which retains a pyramid causeway that is one of the best-preserved in Egypt. Together with the pyramid of Userkaf, this pyramid was the subject of one of the earliest known restoration attempts, conducted by Khaemweset, a son of Ramesses II.[23] Saqqara is also the location of the incomplete step pyramid of Djoser's successor Sekhemkhet, known as the Buried Pyramid. Archaeologists believe that had this pyramid been completed, it would have been larger than Djoser's.
|
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+
|
57 |
+
South of the main pyramid field at Saqqara is a second collection of later, smaller pyramids, including those of Pepi I, Djedkare Isesi, Merenre, Pepi II and Ibi. Most of these are in a poor state of preservation.
|
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|
59 |
+
The Fourth Dynasty pharaoh Shepseskaf either did not share an interest in, or have the capacity to undertake pyramid construction like his predecessors. His tomb, which is also sited at south Saqqara, was instead built as an unusually large mastaba and offering temple complex. It is commonly known as the Mastabat al-Fir’aun.[24]
|
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+
|
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A previously unknown pyramid was discovered at north Saqqara in late 2008. Believed to be the tomb of Teti's mother, it currently stands approximately 5 m (16 ft) high, although the original height was closer to 14 m (46 ft).
|
62 |
+
|
63 |
+
This area is arguably the most important pyramid field in Egypt outside Giza and Saqqara, although until 1996 the site was inaccessible due to its location within a military base and was relatively unknown outside archaeological circles.
|
64 |
+
|
65 |
+
The southern Pyramid of Sneferu, commonly known as the Bent Pyramid, is believed to be the first Egyptian pyramid intended by its builders to be a "true" smooth-sided pyramid from the outset; the earlier pyramid at Meidum had smooth sides in its finished state – but it was conceived and built as a step pyramid, before having its steps filled in and concealed beneath a smooth outer casing of dressed stone. As a true smooth-sided structure, the Bent Pyramid was only a partial success – albeit a unique, visually imposing one; it is also the only major Egyptian pyramid to retain a significant proportion of its original smooth outer limestone casing intact. As such it serves as the best contemporary example of how the ancient Egyptians intended their pyramids to look. Several kilometres to the north of the Bent Pyramid is the last – and most successful – of the three pyramids constructed during the reign of Sneferu; the Red Pyramid is the world's first successfully completed smooth-sided pyramid. The structure is also the third largest pyramid in Egypt – after the pyramids of Khufu and Khafra at Giza.
|
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|
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Also at Dahshur is one of two pyramids built by Amenemhat III, known as the Black Pyramid, as well as a number of small, mostly ruined subsidiary pyramids.
|
68 |
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|
69 |
+
Located to the south of Dahshur, several mudbrick pyramids were built in this area in the late Middle Kingdom, perhaps for Amenemhat IV and Sobekneferu.
|
70 |
+
|
71 |
+
Two major pyramids are known to have been built at Lisht – those of Amenemhat I and his son, Senusret I. The latter is surrounded by the ruins of ten smaller subsidiary pyramids. One of these subsidiary pyramids is known to be that of Amenemhat's cousin, Khaba II.[25] The site which is in the vicinity of the oasis of the Faiyum, midway between Dahshur and Meidum, and about 100 kilometres south of Cairo, is believed to be in the vicinity of the ancient city of Itjtawy (the precise location of which remains unknown), which served as the capital of Egypt during the Twelfth Dynasty.
|
72 |
+
|
73 |
+
The pyramid at Meidum is one of three constructed during the reign of Sneferu, and is believed by some to have been started by that pharaoh's father and predecessor, Huni. However, that attribution is uncertain, as no record of Huni's name has been found at the site. It was constructed as a step pyramid, and then later converted into the first "true" smooth-sided pyramid when the steps were filled in, and an outer casing added. The pyramid suffered several catastrophic collapses in ancient and medieval times; medieval Arab writers described it as having seven steps – although today only the three uppermost of these remain, giving the structure its odd, tower-like appearance. The hill on which the pyramid is situated is not a natural landscape feature – it is the small mountain of debris created when the lower courses and outer casing of the pyramid gave way.
|
74 |
+
|
75 |
+
Amenemhat III was the last powerful ruler of the Twelfth Dynasty, and the pyramid he built at Hawara, near the Faiyum, is believed to post-date the so-called "Black Pyramid" built by the same ruler at Dahshur. It is the Hawara pyramid that is believed to have been Amenemhet's final resting place.
|
76 |
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|
77 |
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The pyramid of Senusret II at el-Lahun is the southernmost royal-tomb pyramid structure in Egypt. Its builders reduced the amount of work necessary to construct it by ingeniously using as its foundation and core a 12-meter-high natural limestone hill.
|
78 |
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|
79 |
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Piye, the king of Kush who became the first ruler of the Twenty-fifth Dynasty of Egypt, built a pyramid at El-Kurru. He was the first Egyptian pharaoh to be buried in a pyramid in centuries.
|
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|
81 |
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Taharqa, a Kushite ruler of the Twenty-fifth Dynasty, built his pyramid at Nuri. It was the largest in the area (North Sudan).
|
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|
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The following table lays out the chronology of the construction of most of the major pyramids mentioned here. Each pyramid is identified through the pharaoh who ordered it built, his approximate reign, and its location.
|
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|
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+
*Would have been 91.65 meters (301 feet) or 175 Egyptian Royal cubits.
|
86 |
+
|
87 |
+
*Originally: 143.5 m or 471 feet or 274 Egyptian Royal cubits
|
88 |
+
|
89 |
+
47.6 m (156 ft; 91 Egyptian Royal cubits)
|
90 |
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|
91 |
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30 meters (99 feet)
|
92 |
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|
93 |
+
50 meters (164 feet)
|
94 |
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|
95 |
+
Constructing the pyramids involved moving huge quantities of stone.
|
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Papyri discovered at the Egyptian desert near the Red Sea, in 2013 by archaeologist Pierre Tallet, revealed the journal of Merer, an official of Egypt involved in transporting limestone along the Nile River. These papyri reveal processes in the building of the Great Pyramid at Giza, the tomb of the Pharaoh Khufu, just outside modern Cairo.[27]
|
97 |
+
Rather than overland transport of the limestone used in building the pyramid, there is evidence that limestone blocks were transported along the Nile River, in the journal of Merer, preserved remnants of ancient canals, and transport boats discovered.[28]
|
98 |
+
|
99 |
+
It is possible that quarried blocks were then transported to the construction site by wooden sleds, with sand in front of the sled wetted to reduce friction. Droplets of water created bridges between the grains of sand, helping them stick together.[29]
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List
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Coordinates: 29°58′34″N 31°07′52″E / 29.97611°N 31.13111°E / 29.97611; 31.13111
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The Egyptian pyramids are ancient pyramid-shaped masonry structures located in Egypt. As of November 2008, sources cite either 118 or 138 as the number of identified Egyptian pyramids.[1][2] Most were built as tombs for the country's pharaohs and their consorts during the Old and Middle Kingdom periods.[3][4][5]
|
6 |
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|
7 |
+
The earliest known Egyptian pyramids are found at Saqqara, northwest of Memphis, although at least one step-pyramid-like structure has been found at Saqqara, dating to the 1st Dynasty: Mastaba 3808, which has been attributed to the reign of Pharaoh Anedjib, with inscriptions, and other archaeological remains of the period, suggesting there may have been others.,[6] the otherwise earliest among these is the Pyramid of Djoser built c. 2630–2610 BC during the Third Dynasty.[7] This pyramid and its surrounding complex are generally considered to be the world's oldest monumental structures constructed of dressed masonry.[8]
|
8 |
+
|
9 |
+
The most famous Egyptian pyramids are those found at Giza, on the outskirts of Cairo. Several of the Giza pyramids are counted among the largest structures ever built.[9] The Pyramid of Khufu at Giza is the largest Egyptian pyramid. It is the only one of the Seven Wonders of the Ancient World still in existence.[10]
|
10 |
+
|
11 |
+
By the time of the Early Dynastic Period, those with sufficient means were buried in bench-like structures known as mastabas.[11][12] At Saqqara, Mastaba 3808, dating from the latter part of the 1st Dynasty, was discovered to contain a large, independently-built step-pyramid-like structure enclosed within the outer palace facade mastaba. Archaeological remains and inscriptions suggest there may have been other similar structures dating to this period.[13]
|
12 |
+
|
13 |
+
The first historically-documented Egyptian pyramid is attributed by Egyptologists to the 3rd Dynasty pharaoh Djoser. Although Egyptologists often credit his vizier Imhotep as its architect, the dynastic Egyptians themselves, contemporaneously or in numerous later dynastic writings about the character, did not credit him with either designing Djoser's pyramid or the invention of stone architecture.[14] The Pyramid of Djoser was first built as a square mastaba-like structure, which as a rule were known to otherwise be rectangular, and was expanded several times by way of a series of accretion layers, to produce the stepped pyramid structure we see today[15] Egyptologists believe this design served as a gigantic stairway by which the soul of the deceased pharaoh could ascend to the heavens.[16]
|
14 |
+
|
15 |
+
Though other pyramids were attempted in the 3rd Dynasty after Djoser, it was the 4th Dynasty, transitioning from the step pyramid to true pyramid shape, which gave rise to the great pyramids of Meidum, Dahsur, and Giza. The last pharaoh of the 4th Dynasty, Shepseskaf, did not build a pyramid and beginning in the 5th Dynasty, for various reasons, the massive scale and precision of construction decreased significantly leaving these later pyramids smaller, less well-built, and often hastily constructed. By the end of the 6th Dynasty pyramid building had largely ended and it was not until the Middle Kingdom that large pyramids were built again, though instead of stone, mud brick was the main construction material.[17]
|
16 |
+
|
17 |
+
Long after the end of Egypt's own pyramid-building period, a burst of pyramid building occurred in what is present-day Sudan, after much of Egypt came under the rule of the Kingdom of Kush, which was then based at Napata. Napatan rule, known as the 25th Dynasty, lasted from 750 BCE to 664 BCE, and during that time Egyptian culture made an indelible impression on the Kushites. The Meroitic period of Kushite history, when the kingdom was centered on Meroë, (approximately in the period between 300 BCE and 300 CE), experienced a full-blown pyramid-building revival, which saw more than two hundred Egyptian-inspired indigenous royal pyramid-tombs constructed in the vicinity of the kingdom's capital cities.
|
18 |
+
|
19 |
+
Al-Aziz Uthman (1171–1198), the second Ayyubid Sultan of Egypt, tried to destroy the Giza pyramid complex. He gave up after only damaging the Pyramid of Menkaure because the task proved too large.[18]
|
20 |
+
|
21 |
+
The shape of Egyptian pyramids is thought to represent the primordial mound from which the Egyptians believed the earth was created. The shape of a pyramid is also thought to be representative of the descending rays of the sun, and most pyramids were faced with polished, highly reflective white limestone, in order to give them a brilliant appearance when viewed from a distance. Pyramids were often also named in ways that referred to solar luminescence. For example, the formal name of the Bent Pyramid at Dahshur was The Southern Shining Pyramid, and that of Senwosret at el-Lahun was Senwosret is Shining.
|
22 |
+
|
23 |
+
While it is generally agreed that pyramids were burial monuments, there is continued disagreement on the particular theological principles that might have given rise to them. One suggestion is that they were designed as a type of "resurrection machine."[19]
|
24 |
+
|
25 |
+
The Egyptians believed the dark area of the night sky around which the stars appear to revolve was the physical gateway into the heavens. One of the narrow shafts that extend from the main burial chamber through the entire body of the Great Pyramid points directly towards the center of this part of the sky. This suggests the pyramid may have been designed to serve as a means to magically launch the deceased pharaoh's soul directly into the abode of the gods.[19]
|
26 |
+
|
27 |
+
All Egyptian pyramids were built on the west bank of the Nile, which, as the site of the setting sun, was associated with the realm of the dead in Egyptian mythology.[20]
|
28 |
+
|
29 |
+
In 1842, Karl Richard Lepsius produced the first modern list of pyramids – now known as the Lepsius list of pyramids – in which he counted 67. A great many more have since been discovered. As of November 2008, 118 Egyptian pyramids have been identified.[3]
|
30 |
+
|
31 |
+
The location of Pyramid 29, which Lepsius called the "Headless Pyramid", was lost for a second time when the structure was buried by desert sands after Lepsius's survey. It was found again only during an archaeological dig conducted in 2008.[21]
|
32 |
+
|
33 |
+
Many pyramids are in a poor state of preservation or buried by desert sands. If visible at all, they may appear as little more than mounds of rubble. As a consequence, archaeologists are continuing to identify and study previously unknown pyramid structures.
|
34 |
+
|
35 |
+
The most recent pyramid to be discovered was that of Sesheshet at Saqqara, mother of the Sixth Dynasty pharaoh Teti, announced on 11 November 2008.[4][22]
|
36 |
+
|
37 |
+
All of Egypt's pyramids, except the small Third Dynasty pyramid of Zawyet el-Amwat (or Zawyet el-Mayitin), are sited on the west bank of the Nile, and most are grouped together in a number of pyramid fields. The most important of these are listed geographically, from north to south, below.
|
38 |
+
|
39 |
+
Abu Rawash is the site of Egypt's most northerly pyramid (other than the ruins of Lepsius pyramid number one)[5]— the mostly ruined Pyramid of Djedefre, son and successor of Khufu. Originally it was thought that this pyramid had never been completed, but the current archaeological consensus is that not only was it completed, but that it was originally about the same size as the Pyramid of Menkaure, which would have placed it among the half-dozen or so largest pyramids in Egypt.
|
40 |
+
|
41 |
+
Its location adjacent to a major crossroads made it an easy source of stone. Quarrying, which began in Roman times, has left little apart from about 15 courses of stone superimposed upon the natural hillock that formed part of the pyramid's core. A small adjacent satellite pyramid is in a better state of preservation.
|
42 |
+
|
43 |
+
Giza is the location of the Pyramid of Khufu (also known as the "Great Pyramid" and the "Pyramid of Cheops"); the somewhat smaller Pyramid of Khafre (or Chephren); the relatively modest-sized Pyramid of Menkaure (or Mykerinus), along with a number of smaller satellite edifices known as "Queen's pyramids"; and the Great Sphinx of Giza.
|
44 |
+
|
45 |
+
Of the three, only Khafre's pyramid retains part of its original polished limestone casing, near its apex. This pyramid appears larger than the adjacent Khufu pyramid by virtue of its more elevated location, and the steeper angle of inclination of its construction – it is, in fact, smaller in both height and volume.
|
46 |
+
|
47 |
+
The Giza pyramid complex has been a popular tourist destination since antiquity and was popularized in Hellenistic times when the Great Pyramid was listed by Antipater of Sidon as one of the Seven Wonders of the Ancient World. Today it is the only one of those wonders still in existence.
|
48 |
+
|
49 |
+
This site, halfway between Giza and Abusir, is the location for two unfinished Old Kingdom pyramids. The northern structure's owner is believed to be pharaoh Nebka, while the southern structure, known as the Layer Pyramid, may be attributable to the Third Dynasty pharaoh Khaba, a close successor of Sekhemkhet. If this attribution is correct, Khaba's short reign could explain the seemingly unfinished state of this step pyramid. Today it stands around 17 m (56 ft) high; had it been completed, it is likely to have exceeded 40 m (130 ft).
|
50 |
+
|
51 |
+
There are a total of fourteen pyramids at this site, which served as the main royal necropolis during the Fifth Dynasty. The quality of construction of the Abusir pyramids is inferior to those of the Fourth Dynasty – perhaps signaling a decrease in royal power or a less vibrant economy. They are smaller than their predecessors, and are built of low-quality local limestone.
|
52 |
+
|
53 |
+
The three major pyramids are those of Niuserre, which is also the best preserved, Neferirkare Kakai and Sahure. The site is also home to the incomplete Pyramid of Neferefre. Most of the major pyramids at Abusir were built using similar construction techniques, comprising a rubble core surrounded by steps of mud bricks with a limestone outer casing. The largest of these Fifth Dynasty pyramids, the Pyramid of Neferirkare Kakai, is believed to have been built originally as a step pyramid some 70 m (230 ft) high and then later transformed into a "true" pyramid by having its steps filled in with loose masonry.
|
54 |
+
|
55 |
+
Major pyramids located here include the Pyramid of Djoser – generally identified as the world's oldest substantial monumental structure to be built of dressed stone – the Pyramid of Userkaf, the Pyramid of Teti and the Pyramid of Merikare, dating to the First Intermediate Period of Egypt. Also at Saqqara is the Pyramid of Unas, which retains a pyramid causeway that is one of the best-preserved in Egypt. Together with the pyramid of Userkaf, this pyramid was the subject of one of the earliest known restoration attempts, conducted by Khaemweset, a son of Ramesses II.[23] Saqqara is also the location of the incomplete step pyramid of Djoser's successor Sekhemkhet, known as the Buried Pyramid. Archaeologists believe that had this pyramid been completed, it would have been larger than Djoser's.
|
56 |
+
|
57 |
+
South of the main pyramid field at Saqqara is a second collection of later, smaller pyramids, including those of Pepi I, Djedkare Isesi, Merenre, Pepi II and Ibi. Most of these are in a poor state of preservation.
|
58 |
+
|
59 |
+
The Fourth Dynasty pharaoh Shepseskaf either did not share an interest in, or have the capacity to undertake pyramid construction like his predecessors. His tomb, which is also sited at south Saqqara, was instead built as an unusually large mastaba and offering temple complex. It is commonly known as the Mastabat al-Fir’aun.[24]
|
60 |
+
|
61 |
+
A previously unknown pyramid was discovered at north Saqqara in late 2008. Believed to be the tomb of Teti's mother, it currently stands approximately 5 m (16 ft) high, although the original height was closer to 14 m (46 ft).
|
62 |
+
|
63 |
+
This area is arguably the most important pyramid field in Egypt outside Giza and Saqqara, although until 1996 the site was inaccessible due to its location within a military base and was relatively unknown outside archaeological circles.
|
64 |
+
|
65 |
+
The southern Pyramid of Sneferu, commonly known as the Bent Pyramid, is believed to be the first Egyptian pyramid intended by its builders to be a "true" smooth-sided pyramid from the outset; the earlier pyramid at Meidum had smooth sides in its finished state – but it was conceived and built as a step pyramid, before having its steps filled in and concealed beneath a smooth outer casing of dressed stone. As a true smooth-sided structure, the Bent Pyramid was only a partial success – albeit a unique, visually imposing one; it is also the only major Egyptian pyramid to retain a significant proportion of its original smooth outer limestone casing intact. As such it serves as the best contemporary example of how the ancient Egyptians intended their pyramids to look. Several kilometres to the north of the Bent Pyramid is the last – and most successful – of the three pyramids constructed during the reign of Sneferu; the Red Pyramid is the world's first successfully completed smooth-sided pyramid. The structure is also the third largest pyramid in Egypt – after the pyramids of Khufu and Khafra at Giza.
|
66 |
+
|
67 |
+
Also at Dahshur is one of two pyramids built by Amenemhat III, known as the Black Pyramid, as well as a number of small, mostly ruined subsidiary pyramids.
|
68 |
+
|
69 |
+
Located to the south of Dahshur, several mudbrick pyramids were built in this area in the late Middle Kingdom, perhaps for Amenemhat IV and Sobekneferu.
|
70 |
+
|
71 |
+
Two major pyramids are known to have been built at Lisht – those of Amenemhat I and his son, Senusret I. The latter is surrounded by the ruins of ten smaller subsidiary pyramids. One of these subsidiary pyramids is known to be that of Amenemhat's cousin, Khaba II.[25] The site which is in the vicinity of the oasis of the Faiyum, midway between Dahshur and Meidum, and about 100 kilometres south of Cairo, is believed to be in the vicinity of the ancient city of Itjtawy (the precise location of which remains unknown), which served as the capital of Egypt during the Twelfth Dynasty.
|
72 |
+
|
73 |
+
The pyramid at Meidum is one of three constructed during the reign of Sneferu, and is believed by some to have been started by that pharaoh's father and predecessor, Huni. However, that attribution is uncertain, as no record of Huni's name has been found at the site. It was constructed as a step pyramid, and then later converted into the first "true" smooth-sided pyramid when the steps were filled in, and an outer casing added. The pyramid suffered several catastrophic collapses in ancient and medieval times; medieval Arab writers described it as having seven steps – although today only the three uppermost of these remain, giving the structure its odd, tower-like appearance. The hill on which the pyramid is situated is not a natural landscape feature – it is the small mountain of debris created when the lower courses and outer casing of the pyramid gave way.
|
74 |
+
|
75 |
+
Amenemhat III was the last powerful ruler of the Twelfth Dynasty, and the pyramid he built at Hawara, near the Faiyum, is believed to post-date the so-called "Black Pyramid" built by the same ruler at Dahshur. It is the Hawara pyramid that is believed to have been Amenemhet's final resting place.
|
76 |
+
|
77 |
+
The pyramid of Senusret II at el-Lahun is the southernmost royal-tomb pyramid structure in Egypt. Its builders reduced the amount of work necessary to construct it by ingeniously using as its foundation and core a 12-meter-high natural limestone hill.
|
78 |
+
|
79 |
+
Piye, the king of Kush who became the first ruler of the Twenty-fifth Dynasty of Egypt, built a pyramid at El-Kurru. He was the first Egyptian pharaoh to be buried in a pyramid in centuries.
|
80 |
+
|
81 |
+
Taharqa, a Kushite ruler of the Twenty-fifth Dynasty, built his pyramid at Nuri. It was the largest in the area (North Sudan).
|
82 |
+
|
83 |
+
The following table lays out the chronology of the construction of most of the major pyramids mentioned here. Each pyramid is identified through the pharaoh who ordered it built, his approximate reign, and its location.
|
84 |
+
|
85 |
+
*Would have been 91.65 meters (301 feet) or 175 Egyptian Royal cubits.
|
86 |
+
|
87 |
+
*Originally: 143.5 m or 471 feet or 274 Egyptian Royal cubits
|
88 |
+
|
89 |
+
47.6 m (156 ft; 91 Egyptian Royal cubits)
|
90 |
+
|
91 |
+
30 meters (99 feet)
|
92 |
+
|
93 |
+
50 meters (164 feet)
|
94 |
+
|
95 |
+
Constructing the pyramids involved moving huge quantities of stone.
|
96 |
+
Papyri discovered at the Egyptian desert near the Red Sea, in 2013 by archaeologist Pierre Tallet, revealed the journal of Merer, an official of Egypt involved in transporting limestone along the Nile River. These papyri reveal processes in the building of the Great Pyramid at Giza, the tomb of the Pharaoh Khufu, just outside modern Cairo.[27]
|
97 |
+
Rather than overland transport of the limestone used in building the pyramid, there is evidence that limestone blocks were transported along the Nile River, in the journal of Merer, preserved remnants of ancient canals, and transport boats discovered.[28]
|
98 |
+
|
99 |
+
It is possible that quarried blocks were then transported to the construction site by wooden sleds, with sand in front of the sled wetted to reduce friction. Droplets of water created bridges between the grains of sand, helping them stick together.[29]
|
100 |
+
|
101 |
+
List
|
102 |
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|
103 |
+
|
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105 |
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|
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|
107 |
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Coordinates: 29°58′34″N 31°07′52″E / 29.97611°N 31.13111°E / 29.97611; 31.13111
|
en/4869.html.txt
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1 |
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4 |
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The Egyptian pyramids are ancient pyramid-shaped masonry structures located in Egypt. As of November 2008, sources cite either 118 or 138 as the number of identified Egyptian pyramids.[1][2] Most were built as tombs for the country's pharaohs and their consorts during the Old and Middle Kingdom periods.[3][4][5]
|
6 |
+
|
7 |
+
The earliest known Egyptian pyramids are found at Saqqara, northwest of Memphis, although at least one step-pyramid-like structure has been found at Saqqara, dating to the 1st Dynasty: Mastaba 3808, which has been attributed to the reign of Pharaoh Anedjib, with inscriptions, and other archaeological remains of the period, suggesting there may have been others.,[6] the otherwise earliest among these is the Pyramid of Djoser built c. 2630–2610 BC during the Third Dynasty.[7] This pyramid and its surrounding complex are generally considered to be the world's oldest monumental structures constructed of dressed masonry.[8]
|
8 |
+
|
9 |
+
The most famous Egyptian pyramids are those found at Giza, on the outskirts of Cairo. Several of the Giza pyramids are counted among the largest structures ever built.[9] The Pyramid of Khufu at Giza is the largest Egyptian pyramid. It is the only one of the Seven Wonders of the Ancient World still in existence.[10]
|
10 |
+
|
11 |
+
By the time of the Early Dynastic Period, those with sufficient means were buried in bench-like structures known as mastabas.[11][12] At Saqqara, Mastaba 3808, dating from the latter part of the 1st Dynasty, was discovered to contain a large, independently-built step-pyramid-like structure enclosed within the outer palace facade mastaba. Archaeological remains and inscriptions suggest there may have been other similar structures dating to this period.[13]
|
12 |
+
|
13 |
+
The first historically-documented Egyptian pyramid is attributed by Egyptologists to the 3rd Dynasty pharaoh Djoser. Although Egyptologists often credit his vizier Imhotep as its architect, the dynastic Egyptians themselves, contemporaneously or in numerous later dynastic writings about the character, did not credit him with either designing Djoser's pyramid or the invention of stone architecture.[14] The Pyramid of Djoser was first built as a square mastaba-like structure, which as a rule were known to otherwise be rectangular, and was expanded several times by way of a series of accretion layers, to produce the stepped pyramid structure we see today[15] Egyptologists believe this design served as a gigantic stairway by which the soul of the deceased pharaoh could ascend to the heavens.[16]
|
14 |
+
|
15 |
+
Though other pyramids were attempted in the 3rd Dynasty after Djoser, it was the 4th Dynasty, transitioning from the step pyramid to true pyramid shape, which gave rise to the great pyramids of Meidum, Dahsur, and Giza. The last pharaoh of the 4th Dynasty, Shepseskaf, did not build a pyramid and beginning in the 5th Dynasty, for various reasons, the massive scale and precision of construction decreased significantly leaving these later pyramids smaller, less well-built, and often hastily constructed. By the end of the 6th Dynasty pyramid building had largely ended and it was not until the Middle Kingdom that large pyramids were built again, though instead of stone, mud brick was the main construction material.[17]
|
16 |
+
|
17 |
+
Long after the end of Egypt's own pyramid-building period, a burst of pyramid building occurred in what is present-day Sudan, after much of Egypt came under the rule of the Kingdom of Kush, which was then based at Napata. Napatan rule, known as the 25th Dynasty, lasted from 750 BCE to 664 BCE, and during that time Egyptian culture made an indelible impression on the Kushites. The Meroitic period of Kushite history, when the kingdom was centered on Meroë, (approximately in the period between 300 BCE and 300 CE), experienced a full-blown pyramid-building revival, which saw more than two hundred Egyptian-inspired indigenous royal pyramid-tombs constructed in the vicinity of the kingdom's capital cities.
|
18 |
+
|
19 |
+
Al-Aziz Uthman (1171–1198), the second Ayyubid Sultan of Egypt, tried to destroy the Giza pyramid complex. He gave up after only damaging the Pyramid of Menkaure because the task proved too large.[18]
|
20 |
+
|
21 |
+
The shape of Egyptian pyramids is thought to represent the primordial mound from which the Egyptians believed the earth was created. The shape of a pyramid is also thought to be representative of the descending rays of the sun, and most pyramids were faced with polished, highly reflective white limestone, in order to give them a brilliant appearance when viewed from a distance. Pyramids were often also named in ways that referred to solar luminescence. For example, the formal name of the Bent Pyramid at Dahshur was The Southern Shining Pyramid, and that of Senwosret at el-Lahun was Senwosret is Shining.
|
22 |
+
|
23 |
+
While it is generally agreed that pyramids were burial monuments, there is continued disagreement on the particular theological principles that might have given rise to them. One suggestion is that they were designed as a type of "resurrection machine."[19]
|
24 |
+
|
25 |
+
The Egyptians believed the dark area of the night sky around which the stars appear to revolve was the physical gateway into the heavens. One of the narrow shafts that extend from the main burial chamber through the entire body of the Great Pyramid points directly towards the center of this part of the sky. This suggests the pyramid may have been designed to serve as a means to magically launch the deceased pharaoh's soul directly into the abode of the gods.[19]
|
26 |
+
|
27 |
+
All Egyptian pyramids were built on the west bank of the Nile, which, as the site of the setting sun, was associated with the realm of the dead in Egyptian mythology.[20]
|
28 |
+
|
29 |
+
In 1842, Karl Richard Lepsius produced the first modern list of pyramids – now known as the Lepsius list of pyramids – in which he counted 67. A great many more have since been discovered. As of November 2008, 118 Egyptian pyramids have been identified.[3]
|
30 |
+
|
31 |
+
The location of Pyramid 29, which Lepsius called the "Headless Pyramid", was lost for a second time when the structure was buried by desert sands after Lepsius's survey. It was found again only during an archaeological dig conducted in 2008.[21]
|
32 |
+
|
33 |
+
Many pyramids are in a poor state of preservation or buried by desert sands. If visible at all, they may appear as little more than mounds of rubble. As a consequence, archaeologists are continuing to identify and study previously unknown pyramid structures.
|
34 |
+
|
35 |
+
The most recent pyramid to be discovered was that of Sesheshet at Saqqara, mother of the Sixth Dynasty pharaoh Teti, announced on 11 November 2008.[4][22]
|
36 |
+
|
37 |
+
All of Egypt's pyramids, except the small Third Dynasty pyramid of Zawyet el-Amwat (or Zawyet el-Mayitin), are sited on the west bank of the Nile, and most are grouped together in a number of pyramid fields. The most important of these are listed geographically, from north to south, below.
|
38 |
+
|
39 |
+
Abu Rawash is the site of Egypt's most northerly pyramid (other than the ruins of Lepsius pyramid number one)[5]— the mostly ruined Pyramid of Djedefre, son and successor of Khufu. Originally it was thought that this pyramid had never been completed, but the current archaeological consensus is that not only was it completed, but that it was originally about the same size as the Pyramid of Menkaure, which would have placed it among the half-dozen or so largest pyramids in Egypt.
|
40 |
+
|
41 |
+
Its location adjacent to a major crossroads made it an easy source of stone. Quarrying, which began in Roman times, has left little apart from about 15 courses of stone superimposed upon the natural hillock that formed part of the pyramid's core. A small adjacent satellite pyramid is in a better state of preservation.
|
42 |
+
|
43 |
+
Giza is the location of the Pyramid of Khufu (also known as the "Great Pyramid" and the "Pyramid of Cheops"); the somewhat smaller Pyramid of Khafre (or Chephren); the relatively modest-sized Pyramid of Menkaure (or Mykerinus), along with a number of smaller satellite edifices known as "Queen's pyramids"; and the Great Sphinx of Giza.
|
44 |
+
|
45 |
+
Of the three, only Khafre's pyramid retains part of its original polished limestone casing, near its apex. This pyramid appears larger than the adjacent Khufu pyramid by virtue of its more elevated location, and the steeper angle of inclination of its construction – it is, in fact, smaller in both height and volume.
|
46 |
+
|
47 |
+
The Giza pyramid complex has been a popular tourist destination since antiquity and was popularized in Hellenistic times when the Great Pyramid was listed by Antipater of Sidon as one of the Seven Wonders of the Ancient World. Today it is the only one of those wonders still in existence.
|
48 |
+
|
49 |
+
This site, halfway between Giza and Abusir, is the location for two unfinished Old Kingdom pyramids. The northern structure's owner is believed to be pharaoh Nebka, while the southern structure, known as the Layer Pyramid, may be attributable to the Third Dynasty pharaoh Khaba, a close successor of Sekhemkhet. If this attribution is correct, Khaba's short reign could explain the seemingly unfinished state of this step pyramid. Today it stands around 17 m (56 ft) high; had it been completed, it is likely to have exceeded 40 m (130 ft).
|
50 |
+
|
51 |
+
There are a total of fourteen pyramids at this site, which served as the main royal necropolis during the Fifth Dynasty. The quality of construction of the Abusir pyramids is inferior to those of the Fourth Dynasty – perhaps signaling a decrease in royal power or a less vibrant economy. They are smaller than their predecessors, and are built of low-quality local limestone.
|
52 |
+
|
53 |
+
The three major pyramids are those of Niuserre, which is also the best preserved, Neferirkare Kakai and Sahure. The site is also home to the incomplete Pyramid of Neferefre. Most of the major pyramids at Abusir were built using similar construction techniques, comprising a rubble core surrounded by steps of mud bricks with a limestone outer casing. The largest of these Fifth Dynasty pyramids, the Pyramid of Neferirkare Kakai, is believed to have been built originally as a step pyramid some 70 m (230 ft) high and then later transformed into a "true" pyramid by having its steps filled in with loose masonry.
|
54 |
+
|
55 |
+
Major pyramids located here include the Pyramid of Djoser – generally identified as the world's oldest substantial monumental structure to be built of dressed stone – the Pyramid of Userkaf, the Pyramid of Teti and the Pyramid of Merikare, dating to the First Intermediate Period of Egypt. Also at Saqqara is the Pyramid of Unas, which retains a pyramid causeway that is one of the best-preserved in Egypt. Together with the pyramid of Userkaf, this pyramid was the subject of one of the earliest known restoration attempts, conducted by Khaemweset, a son of Ramesses II.[23] Saqqara is also the location of the incomplete step pyramid of Djoser's successor Sekhemkhet, known as the Buried Pyramid. Archaeologists believe that had this pyramid been completed, it would have been larger than Djoser's.
|
56 |
+
|
57 |
+
South of the main pyramid field at Saqqara is a second collection of later, smaller pyramids, including those of Pepi I, Djedkare Isesi, Merenre, Pepi II and Ibi. Most of these are in a poor state of preservation.
|
58 |
+
|
59 |
+
The Fourth Dynasty pharaoh Shepseskaf either did not share an interest in, or have the capacity to undertake pyramid construction like his predecessors. His tomb, which is also sited at south Saqqara, was instead built as an unusually large mastaba and offering temple complex. It is commonly known as the Mastabat al-Fir’aun.[24]
|
60 |
+
|
61 |
+
A previously unknown pyramid was discovered at north Saqqara in late 2008. Believed to be the tomb of Teti's mother, it currently stands approximately 5 m (16 ft) high, although the original height was closer to 14 m (46 ft).
|
62 |
+
|
63 |
+
This area is arguably the most important pyramid field in Egypt outside Giza and Saqqara, although until 1996 the site was inaccessible due to its location within a military base and was relatively unknown outside archaeological circles.
|
64 |
+
|
65 |
+
The southern Pyramid of Sneferu, commonly known as the Bent Pyramid, is believed to be the first Egyptian pyramid intended by its builders to be a "true" smooth-sided pyramid from the outset; the earlier pyramid at Meidum had smooth sides in its finished state – but it was conceived and built as a step pyramid, before having its steps filled in and concealed beneath a smooth outer casing of dressed stone. As a true smooth-sided structure, the Bent Pyramid was only a partial success – albeit a unique, visually imposing one; it is also the only major Egyptian pyramid to retain a significant proportion of its original smooth outer limestone casing intact. As such it serves as the best contemporary example of how the ancient Egyptians intended their pyramids to look. Several kilometres to the north of the Bent Pyramid is the last – and most successful – of the three pyramids constructed during the reign of Sneferu; the Red Pyramid is the world's first successfully completed smooth-sided pyramid. The structure is also the third largest pyramid in Egypt – after the pyramids of Khufu and Khafra at Giza.
|
66 |
+
|
67 |
+
Also at Dahshur is one of two pyramids built by Amenemhat III, known as the Black Pyramid, as well as a number of small, mostly ruined subsidiary pyramids.
|
68 |
+
|
69 |
+
Located to the south of Dahshur, several mudbrick pyramids were built in this area in the late Middle Kingdom, perhaps for Amenemhat IV and Sobekneferu.
|
70 |
+
|
71 |
+
Two major pyramids are known to have been built at Lisht – those of Amenemhat I and his son, Senusret I. The latter is surrounded by the ruins of ten smaller subsidiary pyramids. One of these subsidiary pyramids is known to be that of Amenemhat's cousin, Khaba II.[25] The site which is in the vicinity of the oasis of the Faiyum, midway between Dahshur and Meidum, and about 100 kilometres south of Cairo, is believed to be in the vicinity of the ancient city of Itjtawy (the precise location of which remains unknown), which served as the capital of Egypt during the Twelfth Dynasty.
|
72 |
+
|
73 |
+
The pyramid at Meidum is one of three constructed during the reign of Sneferu, and is believed by some to have been started by that pharaoh's father and predecessor, Huni. However, that attribution is uncertain, as no record of Huni's name has been found at the site. It was constructed as a step pyramid, and then later converted into the first "true" smooth-sided pyramid when the steps were filled in, and an outer casing added. The pyramid suffered several catastrophic collapses in ancient and medieval times; medieval Arab writers described it as having seven steps – although today only the three uppermost of these remain, giving the structure its odd, tower-like appearance. The hill on which the pyramid is situated is not a natural landscape feature – it is the small mountain of debris created when the lower courses and outer casing of the pyramid gave way.
|
74 |
+
|
75 |
+
Amenemhat III was the last powerful ruler of the Twelfth Dynasty, and the pyramid he built at Hawara, near the Faiyum, is believed to post-date the so-called "Black Pyramid" built by the same ruler at Dahshur. It is the Hawara pyramid that is believed to have been Amenemhet's final resting place.
|
76 |
+
|
77 |
+
The pyramid of Senusret II at el-Lahun is the southernmost royal-tomb pyramid structure in Egypt. Its builders reduced the amount of work necessary to construct it by ingeniously using as its foundation and core a 12-meter-high natural limestone hill.
|
78 |
+
|
79 |
+
Piye, the king of Kush who became the first ruler of the Twenty-fifth Dynasty of Egypt, built a pyramid at El-Kurru. He was the first Egyptian pharaoh to be buried in a pyramid in centuries.
|
80 |
+
|
81 |
+
Taharqa, a Kushite ruler of the Twenty-fifth Dynasty, built his pyramid at Nuri. It was the largest in the area (North Sudan).
|
82 |
+
|
83 |
+
The following table lays out the chronology of the construction of most of the major pyramids mentioned here. Each pyramid is identified through the pharaoh who ordered it built, his approximate reign, and its location.
|
84 |
+
|
85 |
+
*Would have been 91.65 meters (301 feet) or 175 Egyptian Royal cubits.
|
86 |
+
|
87 |
+
*Originally: 143.5 m or 471 feet or 274 Egyptian Royal cubits
|
88 |
+
|
89 |
+
47.6 m (156 ft; 91 Egyptian Royal cubits)
|
90 |
+
|
91 |
+
30 meters (99 feet)
|
92 |
+
|
93 |
+
50 meters (164 feet)
|
94 |
+
|
95 |
+
Constructing the pyramids involved moving huge quantities of stone.
|
96 |
+
Papyri discovered at the Egyptian desert near the Red Sea, in 2013 by archaeologist Pierre Tallet, revealed the journal of Merer, an official of Egypt involved in transporting limestone along the Nile River. These papyri reveal processes in the building of the Great Pyramid at Giza, the tomb of the Pharaoh Khufu, just outside modern Cairo.[27]
|
97 |
+
Rather than overland transport of the limestone used in building the pyramid, there is evidence that limestone blocks were transported along the Nile River, in the journal of Merer, preserved remnants of ancient canals, and transport boats discovered.[28]
|
98 |
+
|
99 |
+
It is possible that quarried blocks were then transported to the construction site by wooden sleds, with sand in front of the sled wetted to reduce friction. Droplets of water created bridges between the grains of sand, helping them stick together.[29]
|
100 |
+
|
101 |
+
List
|
102 |
+
|
103 |
+
|
104 |
+
|
105 |
+
|
106 |
+
|
107 |
+
Coordinates: 29°58′34″N 31°07′52″E / 29.97611°N 31.13111°E / 29.97611; 31.13111
|