--- license: cc-by-nc-4.0 library_name: diffusers base_model: runwayml/stable-diffusion-v1-5 tags: - lora - text-to-image - diffusers - stable-diffusion inference: False --- # ⚡ Flash Diffusion: FlashSD ⚡ Flash Diffusion is a diffusion distillation method proposed in [ADD ARXIV]() *by Clément Chadebec, Onur Tasar, Eyal Benaroche, and Benjamin Aubin.* This model is a **26.4M** LoRA distilled version of SD1.5 model that is able to generate images in **2-4 steps**. The main purpose of this model is to reproduce the main results of the paper. See our [live demo](https://huggingface.co/spaces/jasperai/FlashPixart).

# How to use? The model can be used using the `StableDiffusionPipeline` from `diffusers` library directly. It can allow reducing the number of required sampling steps to **2-4 steps**. ```python from diffusers import StableDiffusionPipeline, LCMScheduler adapter_id = "jasperai/flash-sd" pipe = StableDiffusionPipeline.from_pretrained( "runwayml/stable-diffusion-v1-5", use_safetensors=True, ) pipe.scheduler = LCMScheduler.from_pretrained( "runwayml/stable-diffusion-v1-5", subfolder="scheduler", timestep_spacing="trailing", ) pipe.to("cuda") # Fuse and load LoRA weights pipe.load_lora_weights(adapter_id) pipe.fuse_lora() prompt = "A raccoon reading a book in a lush forest." image = pipe(prompt, num_inference_steps=4, guidance_scale=0).images[0] ```

# Training Details The model was trained for 20k iterations on 2 H100 GPUs (representing approx. a total **26 GPU hours** of training). Please refer to the [paper]() for further parameters details. **Metrics on COCO 2017 validation set (Table 1)** - FID-5k: 22.6 (2 NFE) / 22.5 (4 NFE) - CLIP Score (ViT-g/14): 0.306 (2 NFE) / 0.311 (4 NFE) **Metrics on COCO 2014 validation (Table 2)** - FID-30k: 12.41 (4 NFE) - FID-30k: 12.27 (2 NFE) ## License This model is released under the the Creative Commons BY-NC license.