license: openrail++ | |
library_name: diffusers | |
tags: | |
- text-to-image | |
- text-to-image | |
- diffusers-training | |
- diffusers | |
- lora | |
- template:sd-lora | |
- stable-diffusion-xl | |
- stable-diffusion-xl-diffusers | |
base_model: stabilityai/stable-diffusion-xl-base-1.0 | |
instance_prompt: A [v28] | |
widget: | |
- text: ' ' | |
output: | |
url: image_0.png | |
<!-- This model card has been generated automatically according to the information the training script had access to. You | |
should probably proofread and complete it, then remove this comment. --> | |
# 'SDXL B-LoRA - lora-library/B-LoRA-drawing3 | |
<Gallery /> | |
## Model description | |
These are lora-library/B-LoRA-drawing3 LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0. | |
The weights were trained using [DreamBooth](https://dreambooth.github.io/). | |
LoRA for the text encoder was enabled: False. | |
Special VAE used for training: madebyollin/sdxl-vae-fp16-fix. | |
## Trigger words | |
You should use "A [v28]" to trigger the image generation. | |
## Download model | |
Weights for this model are available in Safetensors format. | |
[Download](lora-library/B-LoRA-drawing3/tree/main) them in the Files & versions tab. | |
## Intended uses & limitations | |
#### How to use | |
```python | |
# TODO: add an example code snippet for running this diffusion pipeline | |
``` | |
#### Limitations and bias | |
[TODO: provide examples of latent issues and potential remediations] | |
## Training details | |
[TODO: describe the data used to train the model] |