##### `AI_Demo.py` ##### AI Demo, hosted on https://huggingface.co/spaces/DrBenjamin/AI_Demo ##### Please reach out to ben@benbox.org for any questions #### Loading needed Python libraries import streamlit as st import numpy as np import audio2numpy as a2n from pydub import AudioSegment import cv2 from PIL import Image import torch from diffusers import StableDiffusionPipeline from diffusers import StableDiffusionImg2ImgPipeline from diffusers import StableDiffusionControlNetPipeline, ControlNetModel, UniPCMultistepScheduler from transformers import pipeline, set_seed from transformers import VisionEncoderDecoderModel, ViTImageProcessor, AutoTokenizer import os os.environ['COMMANDLINE_ARGS'] = '--skip-torch-cuda-test --precision full --no-half' os.environ['KMP_DUPLICATE_LIB_OK'] = 'True' #### Functions ### Function predict_step = Image to Text recognition def predict_step(image): if image.mode != "RGB": image = image.convert(mode = "RGB") pixel_values = feature_extractor(images = image, return_tensors = "pt").pixel_values pixel_values = pixel_values.to(device) output_ids = model.generate(pixel_values, **gen_kwargs) preds = tokenizer.batch_decode(output_ids, skip_special_tokens = True) preds = [pred.strip() for pred in preds] return str(preds[0]).capitalize() + '.' #### Models st.header('🤗 🤗 Hugging Face Diffusers') st.write('State-of-the-art diffusion models for image, text and audio generation in PyTorch.') devices = ["mps", "cpu", "cuda"] device = st.selectbox(label = 'Select device', options = devices, index = 1, disabled = True) st.write(':orange[MPS for Mac (Metal Performance Shaders), CPU for all systems and CUDA for systems with NVIDIA GPU.]') models = ["runwayml/stable-diffusion-v1-5", "stabilityai/stable-diffusion-2-1", "hakurei/waifu-diffusion", "stabilityai/stable-diffusion-2-base", "nlpconnect/vit-gpt2-image-captioning", "openai-gpt", "gpt2-large", "openai/whisper-large-v2"] model_id_or_path = st.selectbox(label = 'Select model', options = models, index = 5, disabled = True) if model_id_or_path == "runwayml/stable-diffusion-v1-5": st.write(':orange[Stable Diffusion v1-5 is the state of the art text-to-image model.]') elif model_id_or_path == "stabilityai/stable-diffusion-2-1": st.write(':orange[New stable diffusion text-to-image model at 768x768 resolution.]') elif model_id_or_path == "stabilityai/stable-diffusion-2-base": st.write(':orange[New stable diffusion text-to-image model at 512x512 resolution.]') elif model_id_or_path == "hakurei/waifu-diffusion": st.write( ':orange[waifu-diffusion is a latent text-to-image diffusion model that has been conditioned on high-quality anime images through fine-tuning.]') elif model_id_or_path == "nlpconnect/vit-gpt2-image-captioning": st.write(':orange[vit-gpt2 is an image captioning model.]') elif model_id_or_path == "openai-gpt": st.write( ':orange[openai-gpt is a transformer-based language model created and released by OpenAI. The model is a causal (unidirectional) transformer pre-trained using language modeling on a large corpus with long range dependencies.]') elif model_id_or_path == "gpt2-large": st.write( ':orange[GPT-2 Large is the 774M parameter version of GPT-2, a transformer-based language model created and released by OpenAI. The model is a pretrained model on English language using a causal language modeling (CLM) objective.]') elif model_id_or_path == "openai/whisper-large-v2": st.write(':orange[Whisper is a pre-trained model for automatic speech recognition (ASR) and speech translation.]') control_net_models = ["None", "lllyasviel/sd-controlnet-canny", "lllyasviel/sd-controlnet-scribble"] if model_id_or_path == "runwayml/stable-diffusion-v1-5": disable = False else: disable = True control_net_model = st.selectbox(label = 'Select control net model', options = control_net_models, disabled = disable) if control_net_model == "lllyasviel/sd-controlnet-canny": st.write( ':orange[ControlNet is a neural network structure to control diffusion models by adding extra conditions. This checkpoint corresponds to the ControlNet conditioned on Canny edges.]') elif control_net_model == "lllyasviel/sd-controlnet-scribble": st.write( ':orange[ControlNet is a neural network structure to control diffusion models by adding extra conditions. This checkpoint corresponds to the ControlNet conditioned on Scribble images.]') if model_id_or_path != "runwayml/stable-diffusion-v1-5": control_net_model = "None" #### Stable diffusion image 2 image with Control Net if model_id_or_path == "runwayml/stable-diffusion-v1-5" and control_net_model != "None": with st.form('img2img (Control Net)'): st.subheader('Image 2 Image (Control Net)') st.write('Create an image from text input with an image as template.') image = '' uploaded_file = st.file_uploader(label = "Upload a picture", type = 'png') prompt = st.text_input(label = 'Prompt', value = 'A picture in comic style, bright colours, a house with red bricks, a dark sky with a full yellow moon, best quality, extremely detailed.') submitted = st.form_submit_button('Submit') if submitted: # Check for image data if uploaded_file is not None: image = cv2.imdecode(np.frombuffer(uploaded_file.getvalue(), np.uint8), cv2.COLOR_GRAY2BGR) image = cv2.cvtColor(image, cv2.COLOR_BGR2GRAY) # Resize image if existend and not 768x640 / 640x768 pixel h, w = image.shape if not (h == 768 and w == 640) and not (h == 640 and w == 768): # Image is bigger in height than width if h > w: # Resize cropped image to standard dimensions image = cv2.resize(image, (640, 768), interpolation = cv2.INTER_AREA) # Image is smaller in height than width else: # Resize cropped image to standard dimensions image = cv2.resize(image, (768, 640), interpolation = cv2.INTER_AREA) # Get canny image image = cv2.Canny(image, 100, 200) image = image[:, :, None] image = np.concatenate([image, image, image], axis = 2) canny_image = Image.fromarray(image) st.subheader('Preview annotator result') st.image(canny_image) # Load control net and stable diffusion v1-5 controlnet = ControlNetModel.from_pretrained(control_net_model, torch_dtype = torch.float32) pipe = StableDiffusionControlNetPipeline.from_pretrained(model_id_or_path, controlnet = controlnet, torch_dtype = torch.float32) pipe = pipe.to(device) # Recommended if your computer has < 64 GB of RAM pipe.enable_attention_slicing() # Speed up diffusion process with faster scheduler and memory optimization pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config) # Generate image generator = torch.manual_seed(0) image = pipe(prompt = prompt, negative_prompt = "monochrome, lowres, bad anatomy, worst quality, low quality", num_inference_steps = 30, generator = generator, image = canny_image).images[0] st.subheader('Diffuser result') st.write('Model :orange[' + model_id_or_path + '] + :red[' + control_net_model + ']') st.image(image) ## Stable-Diffusion if model_id_or_path == "runwayml/stable-diffusion-v1-5" and control_net_model == "None": with st.form('img2img'): st.subheader('Image 2 Image') st.write('Create an image from text input with an image as template.') image = '' uploaded_file = st.file_uploader(label = "Upload a picture", type = 'png') prompt = st.text_input(label = 'Prompt', value = 'A picture in comic style, bright colours, a house with red bricks, a dark sky with a full yellow moon, best quality, extremely detailed.') submitted = st.form_submit_button('Submit') if submitted: # Check for image data if uploaded_file is not None: image = cv2.imdecode(np.frombuffer(uploaded_file.getvalue(), np.uint8), cv2.IMREAD_COLOR) # Resize image if existend and not 768x640 / 640x768 pixel h, w, _ = image.shape if not (h == 768 and w == 640) and not (h == 640 and w == 768): # Image is bigger in height than width if h > w: # Resize cropped image to standard dimensions image = cv2.resize(image, (640, 768), interpolation = cv2.INTER_AREA) # Image is smaller in height than width else: # Resize cropped image to standard dimensions image = cv2.resize(image, (768, 640), interpolation = cv2.INTER_AREA) image = Image.fromarray(image) # Load the pipeline pipe = StableDiffusionImg2ImgPipeline.from_pretrained(model_id_or_path, torch_dtype = torch.float32) pipe = pipe.to(device) # Recommended if your computer has < 64 GB of RAM pipe.enable_attention_slicing() # Speed up diffusion process with faster scheduler and memory optimization pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config) # Create new image images = pipe(prompt = prompt, negative_prompt = "monochrome, lowres, bad anatomy, worst quality, low quality", num_inference_steps = 30, image = image, strength = 0.75, guidance_scale = 7.5).images # Show image st.subheader('Diffuser result') st.write('Model :orange[' + model_id_or_path + ']') st.image(images[0]) #### Stable diffusion txt 2 image if control_net_model == "None" and model_id_or_path != "nlpconnect/vit-gpt2-image-captioning" and model_id_or_path != "openai-gpt" and model_id_or_path != "gpt2-large" and model_id_or_path != "openai/whisper-large-v2": with st.form('txt2img'): st.subheader('Text 2 Image') st.write('Create an image from text input.') if model_id_or_path == "runwayml/stable-diffusion-v1-5" or model_id_or_path == "stabilityai/stable-diffusion-2-1": value = 'A picture in comic style, bright colours, a house with red bricks, a dark sky with a full yellow moon, best quality, extremely detailed.' if model_id_or_path == "hakurei/waifu-diffusion": value = 'A picture in Anime style, bright colours, a house with red bricks, a dark sky with a full yellow moon, best quality, extremely detailed.' if model_id_or_path == "stabilityai/stable-diffusion-2-base": value = 'A picture in comic style, a castle with grey bricks in the background, a river is going through, a blue sky with a full yellow sun, best quality, extremely detailed.' prompt = st.text_input(label = 'Prompt', value = value) submitted = st.form_submit_button('Submit') if submitted: # Make sure you're logged in with `huggingface-cli login` pipe = StableDiffusionPipeline.from_pretrained(model_id_or_path) pipe = pipe.to(device) # Recommended if your computer has < 64 GB of RAM pipe.enable_attention_slicing() # Speed up diffusion process with faster scheduler and memory optimization pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config) # Results if model_id_or_path == "hakurei/waifu-diffusion": negative = "several scenes, more than one image, split picture" else: negative = "monochrome, lowres, bad anatomy, worst quality, low quality" image = pipe(prompt = prompt, negative_prompt = negative, num_inference_steps = 30, guidance_scale = 7.5).images[0] st.subheader('Diffuser result') st.write('Model :orange[' + model_id_or_path + ']') st.image(image) #### Text (OpenAI gpt models) if model_id_or_path == "openai-gpt" or model_id_or_path == "gpt2-large": with st.form('GPT'): st.subheader('Text generation') st.write('Create text which is generated from text input.') text_input = st.text_input(label = 'Give a start of a sentence', value = 'This is a test ') submitted = st.form_submit_button('Submit') if submitted: generator = pipeline('text-generation', model = model_id_or_path) set_seed(42) generated = generator(text_input, max_length = 150, num_return_sequences = 1) st.subheader('Diffuser result') st.write('Model :orange[' + model_id_or_path + ']') st.markdown('Text: ":green[' + str(generated[0]['generated_text']) + ']"') #### Image to text if model_id_or_path == "nlpconnect/vit-gpt2-image-captioning": with st.form('Image2Text'): st.subheader('Image 2 Text') st.write('Create a description of an image.') image = '' uploaded_file = st.file_uploader(label = "Upload a picture", type = 'png') submitted = st.form_submit_button('Submit') if submitted: # Check for image data if uploaded_file is not None: image = cv2.imdecode(np.frombuffer(uploaded_file.getvalue(), np.uint8), cv2.IMREAD_COLOR) image = Image.fromarray(image) model = VisionEncoderDecoderModel.from_pretrained(model_id_or_path) feature_extractor = ViTImageProcessor.from_pretrained(model_id_or_path) tokenizer = AutoTokenizer.from_pretrained(model_id_or_path) device = torch.device("cuda" if torch.cuda.is_available() else "cpu") model.to(device) max_length = 16 num_beams = 4 gen_kwargs = {"max_length": max_length, "num_beams": num_beams} output = predict_step(image) st.subheader('Diffuser result') st.write('Model :orange[nlpconnect/vit-gpt2-image-captioning]') st.write('Description: ":green[' + str(output) + ']"') #### Whisper Model if model_id_or_path == "openai/whisper-large-v2": with st.form('Image2Text'): st.subheader('Audio 2 Text') st.write('Create a transcription of an audio file.') audio_file = st.file_uploader(label = "Upload an audio file", type = 'mp3') submitted = st.form_submit_button('Submit') if submitted: if audio_file is not None: audio = audio_file.getvalue() with open("temp.mp3", "wb") as binary_file: # Write bytes to file binary_file.write(audio) # Calling the split_to_mono method on the stereo audio file stereo_audio = AudioSegment.from_file("temp.mp3", format = "mp3") mono_audios = stereo_audio.split_to_mono() mono_audios[0].export("temp.mp3", format = "mp3") # Mp3 file to numpy array audio, sr = a2n.audio_from_file('temp.mp3') st.audio('temp.mp3') if os.path.exists("temp.mp3"): os.remove("temp.mp3") # Load model and processor pipe = pipeline("automatic-speech-recognition", model = "openai/whisper-large-v2", chunk_length_s = 30, device = "cpu", ignore_warning = True) prediction = pipe(audio, sampling_rate = sr)["text"] st.subheader('Preview used audio') st.write('Model :orange[' + model_id_or_path + ']') st.write('Transcript: ":green[' + str(prediction) + ']"')