from __future__ import annotations import math import random import spaces import gradio as gr import torch from PIL import Image, ImageOps from diffusers import StableDiffusionInstructPix2PixPipeline help_text = """ If you're not getting what you want, there may be a few reasons: 1. Is the image not changing enough? Your Image CFG weight may be too high. This value dictates how similar the output should be to the input. It's possible your edit requires larger changes from the original image, and your Image CFG weight isn't allowing that. Alternatively, your Text CFG weight may be too low. This value dictates how much to listen to the text instruction. The default Image CFG of 1.5 and Text CFG of 7.5 are a good starting point, but aren't necessarily optimal for each edit. Try: * Decreasing the Image CFG weight, or * Increasing the Text CFG weight, or 2. Conversely, is the image changing too much, such that the details in the original image aren't preserved? Try: * Increasing the Image CFG weight, or * Decreasing the Text CFG weight 3. Try generating results with different random seeds by setting "Randomize Seed" and running generation multiple times. You can also try setting "Randomize CFG" to sample new Text CFG and Image CFG values each time. 4. Rephrasing the instruction sometimes improves results (e.g., "turn him into a dog" vs. "make him a dog" vs. "as a dog"). 5. Increasing the number of steps sometimes improves results. 6. Do faces look weird? The Stable Diffusion autoencoder has a hard time with faces that are small in the image. Try: * Cropping the image so the face takes up a larger portion of the frame. """ example_instructions = [ "Make it a picasso painting", "as if it were by modigliani", "convert to a bronze statue", "Turn it into an anime.", "have it look like a graphic novel", "make him gain weight", "what would he look like bald?", "Have him smile", "Put him in a cocktail party.", "move him at the beach.", "add dramatic lighting", "Convert to black and white", "What if it were snowing?", "Give him a leather jacket", "Turn him into a cyborg!", "make him wear a beanie", ] model_id = "timbrooks/instruct-pix2pix" def main(): pipe = StableDiffusionInstructPix2PixPipeline.from_pretrained(model_id, torch_dtype=torch.float16, safety_checker=None).to("cuda") example_image = Image.open("imgs/example.jpg").convert("RGB") def load_example( steps: int, randomize_seed: bool, seed: int, randomize_cfg: bool, text_cfg_scale: float, image_cfg_scale: float, ): example_instruction = random.choice(example_instructions) return [example_image, example_instruction] + generate( example_image, example_instruction, steps, randomize_seed, seed, randomize_cfg, text_cfg_scale, image_cfg_scale, ) @spaces.GPU(duration=600) def generate( input_image: Image.Image, instruction: str, steps: int, randomize_seed: bool, seed: int, randomize_cfg: bool, text_cfg_scale: float, image_cfg_scale: float, ): seed = random.randint(0, 100000) if randomize_seed else seed text_cfg_scale = round(random.uniform(6.0, 9.0), ndigits=2) if randomize_cfg else text_cfg_scale image_cfg_scale = round(random.uniform(1.2, 1.8), ndigits=2) if randomize_cfg else image_cfg_scale width, height = input_image.size factor = 512 / max(width, height) factor = math.ceil(min(width, height) * factor / 64) * 64 / min(width, height) width = int((width * factor) // 64) * 64 height = int((height * factor) // 64) * 64 input_image = ImageOps.fit(input_image, (width, height), method=Image.Resampling.LANCZOS) if instruction == "": return [input_image, seed] generator = torch.manual_seed(seed) edited_image = pipe( instruction, image=input_image, guidance_scale=text_cfg_scale, image_guidance_scale=image_cfg_scale, num_inference_steps=steps, generator=generator, ).images[0] return [seed, text_cfg_scale, image_cfg_scale, edited_image] def reset(): return [0, "Randomize Seed", 1371, "Fix CFG", 7.5, 1.5, None] with gr.Blocks() as demo: gr.HTML("""
For faster inference without waiting in queue, you may duplicate the space and upgrade to GPU in settings.