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--- |
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license: apache-2.0 |
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datasets: |
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- yuvalkirstain/pickapic_v1 |
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language: |
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- en |
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pipeline_tag: text-to-image |
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--- |
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# Step-aware Preference Optimization: Aligning Preference with Denoising Performance at Each Step |
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<a href="https://arxiv.org/abs/2406.04314"><img src="https://img.shields.io/badge/Paper-arXiv-red?style=for-the-badge" height=22.5></a> |
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<a href="https://github.com/RockeyCoss/SPO"><img src="https://img.shields.io/badge/Gihub-Code-succees?style=for-the-badge&logo=GitHub" height=22.5></a> |
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<a href="https://rockeycoss.github.io/spo.github.io/"><img src="https://img.shields.io/badge/Project-Page-blue?style=for-the-badge" height=22.5></a> |
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<table> |
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<tr> |
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<td><img src="assets/imgs/0.png" alt="teaser example 0" width="200"/></td> |
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<td><img src="assets/imgs/1.png" alt="teaser example 1" width="200"/></td> |
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<td><img src="assets/imgs/2.png" alt="teaser example 2" width="200"/></td> |
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<td><img src="assets/imgs/3.png" alt="teaser example 3" width="200"/></td> |
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</tr> |
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</table> |
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## Abstract |
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<p> |
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Recently, Direct Preference Optimization (DPO) has extended its success from aligning large language models (LLMs) to aligning text-to-image diffusion models with human preferences. |
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Unlike most existing DPO methods that assume all diffusion steps share a consistent preference order with the final generated images, we argue that this assumption neglects step-specific denoising performance and that preference labels should be tailored to each step's contribution. |
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</p> |
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<p> |
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To address this limitation, we propose Step-aware Preference Optimization (SPO), a novel post-training approach that independently evaluates and adjusts the denoising performance at each step, using a <em>step-aware preference model</em> and a <em>step-wise resampler</em> to ensure accurate step-aware supervision. |
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Specifically, at each denoising step, we sample a pool of images, find a suitable win-lose pair, and, most importantly, randomly select a single image from the pool to initialize the next denoising step. This step-wise resampler process ensures the next win-lose image pair comes from the same image, making the win-lose comparison independent of the previous step. To assess the preferences at each step, we train a separate step-aware preference model that can be applied to both noisy and clean images. |
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</p> |
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<p> |
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Our experiments with Stable Diffusion v1.5 and SDXL demonstrate that SPO significantly outperforms the latest Diffusion-DPO in aligning generated images with complex, detailed prompts and enhancing aesthetics, while also achieving more than 20× times faster in training efficiency. Code and model: <a ref="https://rockeycoss.github.io/spo.github.io/">https://rockeycoss.github.io/spo.github.io/</a> |
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</p> |
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## Model Description |
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This model is fine-tuned from [stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0). It has been trained on 4,000 prompts for 10 epochs. This checkpoint is a LoRA checkpoint. We also provide a LoRA checkpoint compatible with [stable-diffusion-webui](https://github.com/AUTOMATIC1111/stable-diffusion-webui), which can be accessed [here](https://civitai.com/models/510261?modelVersionId=567119) |
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If you want to access the merged checkpoint that combines the LoRA checkpoint with the base model [stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0), please visit [SPO-SDXL_4k-p_10ep](https://huggingface.co/SPO-Diffusion-Models/SPO-SDXL_4k-p_10ep). |
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## A quick example |
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```python |
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from diffusers import StableDiffusionXLPipeline, UNet2DConditionModel |
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import torch |
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# load pipeline |
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inference_dtype = torch.float16 |
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pipe = StableDiffusionXLPipeline.from_pretrained( |
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"stabilityai/stable-diffusion-xl-base-1.0", |
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torch_dtype=inference_dtype, |
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) |
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vae = AutoencoderKL.from_pretrained( |
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'madebyollin/sdxl-vae-fp16-fix', |
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torch_dtype=inference_dtype, |
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) |
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pipe.vae = vae |
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pipe.load_lora_weights("SPO-Diffusion-Models/SPO-SDXL_4k-p_10ep_LoRA") |
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pipe.to('cuda') |
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generator=torch.Generator(device='cuda').manual_seed(42) |
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image = pipe( |
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prompt='a child and a penguin sitting in front of the moon', |
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guidance_scale=5.0, |
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generator=generator, |
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output_type='pil', |
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).images[0] |
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image.save('moon.png') |
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``` |
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## Citation |
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If you find our work or codebase useful, please consider giving us a star and citing our work. |
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``` |
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@article{liang2024step, |
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title={Step-aware Preference Optimization: Aligning Preference with Denoising Performance at Each Step}, |
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author={Liang, Zhanhao and Yuan, Yuhui and Gu, Shuyang and Chen, Bohan and Hang, Tiankai and Li, Ji and Zheng, Liang}, |
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journal={arXiv preprint arXiv:2406.04314}, |
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year={2024} |
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} |
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``` |