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--- |
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tags: |
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- stable-diffusion-xl |
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- stable-diffusion-xl-diffusers |
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- text-to-image |
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- diffusers |
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- lora |
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- template:sd-lora |
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widget: |
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- text: 'A photo of <s0><s1> fashion model wearing ' |
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output: |
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url: |
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"image_0.png" |
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- text: 'A photo of <s0><s1> fashion model wearing ' |
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output: |
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url: |
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"image_1.png" |
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- text: 'A photo of <s0><s1> fashion model wearing ' |
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output: |
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url: |
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"image_2.png" |
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- text: 'A photo of <s0><s1> fashion model wearing ' |
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output: |
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url: |
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"image_3.png" |
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base_model: stabilityai/stable-diffusion-xl-base-1.0 |
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instance_prompt: A photo of <s0><s1> fashion model wearing |
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license: openrail++ |
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--- |
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# SDXL LoRA DreamBooth - Usman1921/suit-style-fine-tune-sdxl-lora-10-images |
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<Gallery /> |
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## Model description |
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### These are Usman1921/suit-style-fine-tune-sdxl-lora-10-images LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0. |
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## Download model |
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### Use it with UIs such as AUTOMATIC1111, Comfy UI, SD.Next, Invoke |
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- **LoRA**: download **[`suit-style-fine-tune-sdxl-lora-10-images.safetensors` here 💾](/Usman1921/suit-style-fine-tune-sdxl-lora-10-images/blob/main/suit-style-fine-tune-sdxl-lora-10-images.safetensors)**. |
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- Place it on your `models/Lora` folder. |
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- On AUTOMATIC1111, load the LoRA by adding `<lora:suit-style-fine-tune-sdxl-lora-10-images:1>` to your prompt. On ComfyUI just [load it as a regular LoRA](https://comfyanonymous.github.io/ComfyUI_examples/lora/). |
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- *Embeddings*: download **[`suit-style-fine-tune-sdxl-lora-10-images_emb.safetensors` here 💾](/Usman1921/suit-style-fine-tune-sdxl-lora-10-images/blob/main/suit-style-fine-tune-sdxl-lora-10-images_emb.safetensors)**. |
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- Place it on it on your `embeddings` folder |
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- Use it by adding `suit-style-fine-tune-sdxl-lora-10-images_emb` to your prompt. For example, `A photo of suit-style-fine-tune-sdxl-lora-10-images_emb fashion model wearing` |
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(you need both the LoRA and the embeddings as they were trained together for this LoRA) |
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## Use it with the [🧨 diffusers library](https://github.com/huggingface/diffusers) |
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```py |
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from diffusers import AutoPipelineForText2Image |
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import torch |
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from huggingface_hub import hf_hub_download |
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from safetensors.torch import load_file |
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pipeline = AutoPipelineForText2Image.from_pretrained('stabilityai/stable-diffusion-xl-base-1.0', torch_dtype=torch.float16).to('cuda') |
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pipeline.load_lora_weights('Usman1921/suit-style-fine-tune-sdxl-lora-10-images', weight_name='pytorch_lora_weights.safetensors') |
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embedding_path = hf_hub_download(repo_id='Usman1921/suit-style-fine-tune-sdxl-lora-10-images', filename='suit-style-fine-tune-sdxl-lora-10-images_emb.safetensors', repo_type="model") |
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state_dict = load_file(embedding_path) |
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pipeline.load_textual_inversion(state_dict["clip_l"], token=["<s0>", "<s1>"], text_encoder=pipeline.text_encoder, tokenizer=pipeline.tokenizer) |
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pipeline.load_textual_inversion(state_dict["clip_g"], token=["<s0>", "<s1>"], text_encoder=pipeline.text_encoder_2, tokenizer=pipeline.tokenizer_2) |
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image = pipeline('A photo of <s0><s1> fashion model wearing ').images[0] |
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``` |
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For more details, including weighting, merging and fusing LoRAs, check the [documentation on loading LoRAs in diffusers](https://huggingface.co/docs/diffusers/main/en/using-diffusers/loading_adapters) |
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## Trigger words |
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To trigger image generation of trained concept(or concepts) replace each concept identifier in you prompt with the new inserted tokens: |
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to trigger concept `TOK` → use `<s0><s1>` in your prompt |
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## Details |
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All [Files & versions](/Usman1921/suit-style-fine-tune-sdxl-lora-10-images/tree/main). |
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The weights were trained using [🧨 diffusers Advanced Dreambooth Training Script](https://github.com/huggingface/diffusers/blob/main/examples/advanced_diffusion_training/train_dreambooth_lora_sdxl_advanced.py). |
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LoRA for the text encoder was enabled. False. |
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Pivotal tuning was enabled: True. |
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Special VAE used for training: madebyollin/sdxl-vae-fp16-fix. |
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