| library_name: diffusers | |
| license: mit | |
| ## Juggernaut Diffusion XL Model Card | |
| Juggernaut Diffusion XL is a latent text-to-image diffusion model capable of generating images of people, mainly, and other things. For more information about how Stable Diffusion functions, please have a look at 🤗's [Stable Diffusion blog](https://huggingface.co/blog/stable_diffusion). | |
| You can use this with the 🧨Diffusers library from [Hugging Face](https://huggingface.co). | |
|  | |
| ### Diffusers | |
| ```py | |
| from diffusers import StableDiffusionPipeline | |
| import torch | |
| pipeline = StableDiffusionPipeline.from_pretrained("nroggendorff/juggernautxl").to("cuda") | |
| image = pipeline(prompt="a girl standing in the rain").images[0] | |
| image.save("not.png") | |
| ``` | |
| ### Model Details | |
| - `train_batch_size`: 1 | |
| - `gradient_accumulation_steps`: 4 | |
| - `learning_rate`: 1e-2 | |
| - `lr_warmup_steps`: 500 | |
| - `mixed_precision`: "fp16" | |
| - `eval_metric`: "mean_squared_error" | |
| ### Limitations | |
| - The model does not achieve perfect photorealism | |
| - The model cannot render legible text | |
| ### Developed by | |
| - Noa Linden Roggendorff | |
| *This model card was written by Noa Roggendorff and is based on the [Stable Diffusion v1-5 Model Card](https://huggingface.co/runwayml/stable-diffusion-v1-5).* |