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--- |
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license: cc-by-nc-4.0 |
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library_name: diffusers |
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base_model: PixArt-alpha/PixArt-XL-2-1024-MS |
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tags: |
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- lora |
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- text-to-image |
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inference: False |
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--- |
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# ⚡ Flash Diffusion: FlashPixart ⚡ |
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Flash Diffusion is a diffusion distillation method proposed in [Flash Diffusion: Accelerating Any Conditional |
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Diffusion Model for Few Steps Image Generation](http://arxiv.org/abs/2406.02347) *by Clément Chadebec, Onur Tasar, Eyal Benaroche, and Benjamin Aubin* from Jasper Research. |
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This model is a **66.5M** LoRA distilled version of [Pixart-α](https://huggingface.co/PixArt-alpha/PixArt-XL-2-1024-MS) model that is able to generate 1024x1024 images in **4 steps**. |
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See our [live demo](https://huggingface.co/spaces/jasperai/FlashPixart) and official [Github repo](https://github.com/gojasper/flash-diffusion). |
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<p align="center"> |
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<img style="width:700px;" src="assets/flash_pixart.jpg"> |
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</p> |
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# How to use? |
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The model can be used using the `PixArtAlphaPipeline` from `diffusers` library directly. It can allow reducing the number of required sampling steps to **4 steps**. |
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```python |
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import torch |
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from diffusers import PixArtAlphaPipeline, Transformer2DModel, LCMScheduler |
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from peft import PeftModel |
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# Load LoRA |
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transformer = Transformer2DModel.from_pretrained( |
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"PixArt-alpha/PixArt-XL-2-1024-MS", |
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subfolder="transformer", |
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torch_dtype=torch.float16 |
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) |
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transformer = PeftModel.from_pretrained( |
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transformer, |
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"jasperai/flash-pixart" |
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) |
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# Pipeline |
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pipe = PixArtAlphaPipeline.from_pretrained( |
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"PixArt-alpha/PixArt-XL-2-1024-MS", |
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transformer=transformer, |
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torch_dtype=torch.float16 |
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) |
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# Scheduler |
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pipe.scheduler = LCMScheduler.from_pretrained( |
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"PixArt-alpha/PixArt-XL-2-1024-MS", |
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subfolder="scheduler", |
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timestep_spacing="trailing", |
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) |
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pipe.to("cuda") |
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prompt = "A raccoon reading a book in a lush forest." |
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image = pipe(prompt, num_inference_steps=4, guidance_scale=0).images[0] |
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``` |
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<p align="center"> |
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<img style="width:400px;" src="assets/raccoon.png"> |
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</p> |
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# Training Details |
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The model was trained for 40k iterations on 4 H100 GPUs (representing approximately 188 hours of training). Please refer to the [paper](http://arxiv.org/abs/2406.02347) for further parameters details. |
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**Metrics on COCO 2014 validation (Table 4)** |
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- FID-10k: 29.30 (4 NFE) |
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- CLIP Score: 0.303 (4 NFE) |
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## Citation |
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If you find this work useful or use it in your research, please consider citing us |
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```bibtex |
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@misc{chadebec2024flash, |
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title={Flash Diffusion: Accelerating Any Conditional Diffusion Model for Few Steps Image Generation}, |
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author={Clement Chadebec and Onur Tasar and Eyal Benaroche and Benjamin Aubin}, |
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year={2024}, |
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eprint={2406.02347}, |
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archivePrefix={arXiv}, |
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primaryClass={cs.CV} |
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} |
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``` |
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## License |
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This model is released under the the Creative Commons BY-NC license. |