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--- |
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license: creativeml-openrail-m |
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base_model: stabilityai/stable-diffusion-2-1-base |
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instance_prompt: sksdog |
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tags: |
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- stable-diffusion |
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- stable-diffusion-diffusers |
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- text-to-image |
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- diffusers |
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inference: true |
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--- |
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# Openjourney LoRA - by [PromptHero](https://prompthero.com/?utm_source=huggingface&utm_medium=referral) |
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These are LoRA adaption weights for [Openjourney](https://huggingface.co/prompthero/openjourney) trained by [@JHawkk](https://prompthero.com/JHawkk) |
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# Openjourney Links |
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- [Openjourney Dreambooth](https://huggingface.co/prompthero/openjourney) |
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- [Openjourney Fine tuned model](https://huggingface.co/prompthero/openjourney-v2) |
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# Want to learn AI art generation?: |
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- [Crash course in AI art generation](https://prompthero.com/academy/prompt-engineering-course?utm_source=huggingface&utm_medium=referral) |
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- [Learn to fine-tune Stable Diffusion for photorealism](https://prompthero.com/academy/dreambooth-stable-diffusion-train-fine-tune-course?utm_source=huggingface&utm_medium=referral) |
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# How to use LoRA's in auto1111: |
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- Update webui (use git pull like here or redownload it) |
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- Copy the file to stable-diffusion-webui/models/lora |
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- Select your LoRA like in this video |
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- Make sure to change the weight (by default it's :1 which is usually too high) |
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# Examples: |
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